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partial equations

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partial equations

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(Theory)

Claudio Canuto

Dipartimento di Matematica

Politecnico di Torino

10129 Torino Italy

claudio.canuto@polito.it

http://calvino.polito.it/ccanuto

2

Chapter 1

Basic Concepts

1.1 Vectors

The inner product between two column vectors a = (ai ) IRm and b = (bi ) IRm will be denoted

by

m

X

a b := ai bi .

i=1

matrix), since a b is the matrix product between the (1 m)-matrix aT and the (m, 1)-matrix b.

For m = 3, the vector product (or external product) between a and b is the vector

e1 e2 e2

a b := det a1 a2 a3

b1 b2 b3

= (a2 b3 a3 b2 )e1 + (a3 b1 a1 b3 )e2 + (a1 b2 a2 b1 )e3

= (a2 b3 a3 b2 , a3 b1 a1 b3 , a1 b2 a2 b1 )T ,

Let us denote by x = (x1 , . . . , xm ) the independent variable in the Euclidean space IRm . Let

u = u(x ) be a real-valued function defined in some open set O IRm . Given a multi-integer =

(1 , . . . , m ) INm , the -partial derivative of u at a point x O is obtained by differentiating u

at x i -times with respect to the variable xi , for i = 1, . . . , m. The order of the partial derivative

is defined as || := 1 + + m . The -partial derivative will be denoted by one of the symbols

|| u

D u or .

x1 1 . . . xmm

Other symbols may be preferred for indicating low order derivatives. For instance, the first order

partial derivatives with respect to xi will also be denoted by

Di u or Dxi u or ux i ;

3

4 CHAPTER 1. BASIC CONCEPTS

2

Dij u or Dx2i xj u or ux i x j ,

and so on.

We now introduce the most commonly used first order differential operators. The gradient

or grad is defined as

u

x1 x1

. .

u :=

.. = .. u ;

u

xm xm

note that acts on a scalar function and produces a column vector function, i. e., a vector field

defined in O.

The divergence or div acts on a vector-valued function u = (u1 , . . . , um )T and produces a

scalar function, according to the definition

u1 um

u := + + .

x1 xm

The notation is coherent with the fact that u can be formally obtained as the inner product

of the column vectors and u. Therefore, an equivalent notation for u is T u.

In dimension m = 3, the curl or rot acts on a vector-valued function u and produces a

vector-valued function, according to the definition

u3 u2 u1 u3 u2 u1 T

u= , , ;

x2 x3 x3 x1 x1 x2

the vector u can be formally obtained as the vector product of the column vectors and u.

In dimension m = 2, we can define the curl of a scalar function u as the column vector in IR2

u u T

u= , ;

x2 x1

note that u contains the two first components of the vector U IR3 , where U =

(0, 0, u)T IR3 (the last component is obviously 0). Similarly, we can define the curl of a vector

function u = (u1 , u2 )T IR2 as the scalar

u2 u1

u= ,

x1 x2

which coincides with the third component of the vector U IR3 , where U = (u, 0)T (the first

two components are 0 since u does not depend on x3 ).

The perhaps most popular second order differential operator is the Laplacian , defined as

2u 2u

u := + + .

x21 x2m

for this reason, the Laplacian is also denoted by the symbol 2 .

1.2. INTRODUCTION AND NOTATIONS 5

R(x , u; D ) = 0 (1.2.1)

among the independent variable x , the dependent variable u = u(x ) and certain partial derivatives

D applied to u or to some functions depending on u; the multi-integers vary in some finite

subset of INm . The equation is required to be satisfied in some open set O IRm . The order of

the equation is the maximum order of the partial derivatives which appear in the relationship.

Examples of (first order) partial differential equations are

m

X u

a u = ai = 0, (1.2.2)

xi

i=1

u u2

+ = 0, (1.2.3)

t x 2

where we have set (x1 , x2 ) = (x, t) IR2 ;

m

X

u 2

u u = = f, (1.2.4)

xi

i=1

L(x , u; D ) = f, (1.2.5)

and f = f (x ) is a given function in O. Equivalently, the left-hand side of (1.2.5) is a linear

differential operator (of some order N ) applied to u; this means that for all with || N there

exist coefficients a = a(x ) such that

X

L(x , u; D ) = aD u, (1.2.6)

||N

and, at each x O, at least one coefficient with || = N is not vanishing. For convenience, the

left-hand side L(x , u; D ) will be denoted by Lu.

If we restrict the sum in (1.2.6) to the indices with || = N , we obtain the principal part L(N )

of the operator L, i.e., X

L(N ) u = aD u.

||=N

The transport equation (1.2.2) is an example of a linear first order equation. Examples of

linear second order equations (in two independent variables) are

6 CHAPTER 1. BASIC CONCEPTS

2u 2u

+ 2 =f (1.2.7)

x2 y

(called the Laplace equation if f = 0);

u 2 u

2 =f (1.2.8)

t x

(t denotes the time variable, whereas x denotes the space variable);

2u 2u

2 = f; (1.2.9)

t2 x

(iv) the Tricomi equation

2u 2u

+ y = f. (1.2.10)

x2 y 2

Note that the principal part of the heat operator is minus the second order derivative in the x

space variable, i.e., minus the Laplacian in one space variable.

if it can be written as X

aDu = f, (1.2.11)

||=N

where the coefficients a as well as f may depend not only on x but also on u and certain

derivatives D u of order || < N . An example is the inviscid Burgers equation (1.2.3), which can

be written in the (formally) equivalent expression

u u

+u = 0.

t x

Finally, a partial differential equation is semi-linear if it is quasi-linear and the coefficients

a in (1.2.11) depend neither on u nor on its derivatives (whereas f may depend). Examples of

semi-linear equations are the viscous Burgers equation

u 2u u2

2 + = 0,

t x x 2

u u 3 u

+u + = 0, (1.2.12)

t x x3

and the ground-state equation

u = u3 .

We will now discuss in which sense a function u defined in the open set O IRm is a solution

of the partial differential equation (1.2.1) therein. Indeed, we can give different meanings to the

word solution. We go from the concept of classical solution to that of strong solution, and

then to weaker and weaker definitions, which require a solution to be less and less regular (i.e.,

1.2. INTRODUCTION AND NOTATIONS 7

differentiable). One of the main achievements of the Mathematics of the XXth century has been

the relaxation of the concept of solution of a partial differential equation; this has allowed the

differential problems to be formulated in the most appropriate way for being studied by often

sophisticated analytical tools, and numerically discretized by efficient methods.

Let us denote by N the order of the partial differential equation (1.2.1). A classical solution is

a N -time continuously differentiable function in O (i.e., u CN (O)) which, inserted with all its

derivatives in the left-hand side of (1.2.1), makes the equation satisfied pointwise in O:

R(x , u; D ) = 0, x O. (1.2.13)

We want to formulate these conditions in an equivalent way, which subsequently will allow us

to relax the concept of solution. To this end, we introduce the notion of test function, i.e., an

infinitely differentiable function defined and having compact support in O: this means that the

closed set

supp = closure of {x O : (x ) 6= 0}

is bounded and contained in O. Then, vanishes with all its derivatives in a neighborhood of the

boundary O. The set of all test functions forms a vector space, which will be denoted by D(O).

Note that any partial derivative of a test function is itself a test function.

Example 1.2.1. It is often important to know that test functions with certain properties exist;

for example, one often needs a test function that is positive in a small neighborhood of a given

point x 0 and zero outside that neighborhood. Such a function can be given explicitly:

exp 2

if kx x 0 k <

(x ) = kx x 0 k2 2

0 otherwise,

Let us assume that R depends continuously on all its arguments, so that R(x , v; D ) is a

continuous function in O, for all functions v CN (O). The set of conditions (1.2.13) is equivalent

to the set of conditions

Z

R(x , u; D )(x ) dx = 0, D(O). (1.2.14)

O

both sides by (x ) and integrate over O to get (1.2.14). Conversely, suppose that (1.2.14) holds;

assume by contradiction that there exists x 0 O such that R(x 0 , u; D ) 6= 0, say, strictly

positive. Since R(x , u; D ) is a continuous function of x , it will also be strictly positive in

a neighborhood B(x 0 ) of x 0 . Take as test function a nonnegative function having support

contained in B(x 0 ) and satisfying (x 0 ) = 1. Then,

Z

R(x , u; D )(x ) dx > 0,

O

The interest of formulation (1.2.14) relies on the fact that certain derivatives applied to u, or

to some functions of u, can be moved on , thus relaxing the differentiability requirements on u.

This is accomplished via the (repeated) use of the integration-by-parts formula

Z Z

g

(x )(x ) dx = g(x ) (x ) dx , D(O), i = 1, . . . , m, (1.2.15)

O xi O xi

8 CHAPTER 1. BASIC CONCEPTS

which holds, at least, if g is continuously differentiable in O. Note that no boundary term appears,

since a test function vanishes in a neighborhood of O. While the left-hand side requires the partial

derivative of g with respect to xi to be defined and integrable on O, the right-hand side is defined

under the milder condition that g be integrable on O, only.

To explain how (1.2.15) is used to manipulate (1.2.14), assume that the partial differential

equation is written in the quasi-divergence form

m

X

R(x , u; D ) = g(x , u; D ) + g0 (x , u; D ) = gi (x , u; D ) + g0 (x , u; D ),

xi

i=1

where each gi (i = 0, 1, . . . , m) only involves partial derivatives of order strictly less than N .

Many partial differential equations which model fundamental phenomena of the physical world

are precisely obtained in this form; conservation laws are an example. Then, applying (1.2.15),

conditions (1.2.14) can be written as

Z "X m

#

gi (x , u; D ) (x ) + g0 (x , u; D )(x ) dx = 0, D(O). (1.2.16)

O xi

i=1

In this formulation, u need not be differentiable up to order N ; it is enough for the functions gi to

be defined and integrable on O. Any function u for which this is true and which satisfies (1.2.16)

is called a weak solution of the partial differential equation. Obviously, a classical solution is also

a weak solution, whereas the converse need not be true.

Further integrations by parts in (1.2.16) may lead to even weaker definitions of solution.

Example 1.2.2. Consider the transport equation (1.2.2) and assume that the coefficients ai

(i = 1, . . . , m) belong to C1 (O). After a change of sign, the equation can be written as

m m

!

X X ai

(ai u) + u = 0.

xi xi

i=1 i=1

Z "m m

! #

X X ai

u(x ) ai (x ) (x ) + (x ) (x ) dx = 0, D(O).

O xi xi

i=1 i=1

In this way, we allow the transport equation to have bounded, piecewise smooth but discontinuous

weak solutions.

Example 1.2.3. Recalling that = , the Poisson equation

u = f

in O is written in weak form as

Z Z Xm Z

u

u dx = dx = f dx , D(O),

O O xi xi O

i=1

provided f and all the first order partial derivatives of u exist and are integrable on O. A further

integration by parts yields

Z Z

u dx = f dx , D(O),

O O

1.2. INTRODUCTION AND NOTATIONS 9

Partial differential equations are usually supplemented by boundary and/or initial conditions,

i.e., conditions that the solution has to satisfy on all or part of the boundary O of the region O in

which the equation is set; if O is unbounded, the solution may be required to match a prescribed

asymptotic behaviour at infinity. Indeed, in most cases, a partial differential equation admits

infinitely many solutions; the conditions on O or at infinity, which often originate as part of

the mathematical model describing the phenomenon of interest, allow us to select precisely one

solution.

Example 1.2.4. Consider the simple transport equation in one space variable

ut + ux .

It is immediate to check that if g = g(s) is any continuously differentiable function on the real

line, then u(x, t) = g(x t) is a classical solution of the equation. Thus, if we set the equation in

the half-plane {(x, t) IR2 : t > 0}, so that O = {(x, 0) : x IR} represents the space at the

initial time t = 0, then u is the unique solution of the initial value problem

ut + ux = 0 in O

u = g on O.

u = f

in some O IRm . If we know one solution uf , then all solutions can be written in the form

u = u0 + uf , where u0 denotes any harmonic function in O, i.e., any solution of the homogeneous

equation u0 = 0 therein. We shall see that, under appropriate assumptions, a unique solution

can be selected by forcing u to vanish on the whole of O, and, if O is unbounded, at infinity.

The two previous examples concern linear partial differential equations,

Lu = f in O. (1.2.17)

For such equations, the set of solutions is an (affine) vector space. To see this, consider at first

the associated homogeneous equation

Lu0 = 0 in O.

The set of its solutions is a linear vector space: indeed, if u0 and v0 are two such solutions, then

by the linearity of L one has

L(u0 + v0 ) = Lu0 + Lv0 = 0, , IR,

so u0 + v0 is also a solution. Going back to the nonhomogeneous equation, if u and v are two

solutions, then

L(u v) = Lu Lv = f f = 0,

i.e., u v is a solution of the homogeneous equation. Thus, if (1.2.17) admits a solution uf , then

all its solutions can be written in the form

u = u 0 + uf ,

with u0 arbitrary solution of the homogeneous equation. This is the well-known superposition

principle of linear equations.

10 CHAPTER 1. BASIC CONCEPTS

The most general linear first order partial differential equation is

m

X u

a u + a0 u = ai + a0 u = f ; (1.3.1)

xi

i=1

a = (a1 , . . . , am )T 6= 0 and a0 are the coefficients of the equation, whereas f is a given function.

An alternative formulation of the equation is the quasi-divergence form

m

X

(au) + a0 u = (ai u) + a0 u = f.

xi

i=1

If the coefficients ai (i = 1, . . . , m) are differentiable, the two formulations are equivalent by the

differentiation rule of a product, up to a different definition of the zeroth-order coefficient a0 .

We want to show that (1.3.1) is equivalent to a family of ordinary differential equations. We

u

write a = kak a, with a having unitary Euclidean norm, and we denote by = a u the

a

directional derivative of u along a. Then, (1.3.1) becomes

u

kak + a0 u = f,

a

which is a family of ordinary differential equations in the directions of a. To be more explicit, let

us assume that the coefficients ai (i = 1, . . . , m) are bounded, continuously differentiable functions

in the closure O of the region O in which the equation is set. Let us introduce the characteristics

curves of the equation, i.e., the curves x = x (s) defined as the solutions of the autonomous

ordinary differential system

dx

= a(x ). (1.3.2)

ds

Here, s is a real variable which parametrizes each curve. A classical result in the theory of

ordinary differential equations (see, e.g., (??)) guarantees that, under the assumptions made on

the coefficients, for each x 0 O there exists exactly one characteristic curve passing through x 0 ;

it is defined as the solution x = x (s; x 0 ) of the Cauchy problem

dx = a(x )

ds

x (0) = x 0 .

The solution exists for positive and negative values of the parameter s, until x reaches the boundary

O. Note that the characteristics only depend on the principal part of the operator.

Let now u be a (classical) solution of (1.3.1), and let us consider its restriction u = u(s) =

u(x (s)) to a characteristic curve. By the chain rule and (1.3.2), one has

m

du X u dxi

= = a u.

ds xi ds

i=1

It follows that u can be determined by solving the linear ordinary differential equation

du

+ a0 u = f (1.3.3)

ds

1.3. LINEAR FIRST ORDER EQUATIONS 11

O0 n

a

n a n a

O O +

O0 a

n

Figure 1.1: Decomposition of the boundary of a channel O into inflow boundary O , character-

istic boundary O0 and outflow boundary O+

on each characteristic curve (again, the symbol indicates restriction to the characteristic curve);

solvability is guaranteed if, for instance, a0 and f are bounded and continuous in O. Furthermore,

u can be uniquely determined by prescribing its value at one point of each characteristic curve.

A situation of particular interest is the following one. Let O be smooth enough so that the

unit vector n = n(x ) normal to O exists at each point x O; we assume that O is locally on

one side of O, and n is pointing outwards. Let us introduce the inflow boundary of O as the set

The terminology comes from the fact that if a is the (Eulerian) velocity of fluid particles, then O

is the portion of the boundary where the fluid is entering the region O. The sets O+ (outflow

boundary) and O0 (characteristic boundary) are defined similarly, with < replaced by > and

=, respectively.

Now, suppose that each point in O is reached by a characteristic curve issuing from O (see

Figure 1.1). Then, we can prescribe the value of u at each point in O and uniquely solve the

set of equations (1.3.3), getting u at each point in O. In other words, given a function g on O ,

the boundary value problem

a u + a0 u = f in O

(1.3.5)

u = g on O

admits a unique solution.

Before presenting an example, we anticipate that in Chapter ?? we shall see that this problem

is indeed solvable under weaker assumptions on the data (the domain, the coefficients of the

operator and the right-hand sides f and g).

Example 1.3.1. Consider the simple, constant coefficient equation

ut + aux = 0 (1.3.6)

in the variables (x1 , x2 ) = (x, t). Thus, a = (a, 1)T and a0 = 0. The characteristic curves are

defined by the relations

dx dt

= a, = 1.

ds ds

Eliminating s, we get

x at = constant; (1.3.7)

12 CHAPTER 1. BASIC CONCEPTS

t = t0 + a1 (x x0 )

t0 + t

t0

x0 x0 + at x

in other words, the characteristics are straight lines in the plane (x, t) having slope 1/a (see Fig.

1.2). The solution u is constant along these lines. Thus, the equation models the propagation of

a signal in the x-direction, with speed a: a signal issued at time t0 from position x0 is received

at time t0 + t at position x0 + at (see Fig. 1.3). Indeed,

At first, let us suppose that O is the half-plane {(x, t) : t > 0}. Since n = (0, 1)T on

O = {(x, 0) : x IR}, we have a n = 1 therein, so that O = O, i.e., all the boundary is

inflow. Thus, we prescribe the value u0 of u on O, i.e., at the initial time t = 0. In this case,

(1.3.5) reads as

ut + aux = 0 x IR, t > 0 ,

(1.3.8)

u(x, 0) = u0 (x) x IR ,

which is more properly called an initial value problem. Given a point (x, t) O, the characteristic

line passing through it originates from the point (x0 , 0) O such that x at = x0 (see (1.3.7)

and Fig. 1.4). Since u is constant on this line, we have u(x, t) = u(x0 , 0) = u0 (x0 ) = u0 (x at).

1.4 1.4

1.2 1.2

S S

1 1

0.8 0.8

0.4 0.4

0.2 0.2

x x

0 0

x0 x0 +at

0.2 0.2

0 0.2 0.4 0.6 0.8 1 1.2 1.4 1.6 1.8 2 0 0.2 0.4 0.6 0.8 1 1.2 1.4 1.6 1.8 2

1.3. LINEAR FIRST ORDER EQUATIONS 13

(x, t)

x at x

Dropping the bars on x and t, we get the explicit formula for the solution of the initial value

problem (1.3.8)

u(x, t) = u0 (x at) , for all (x, t) O . (1.3.9)

with respect to x and t; thus, u is a classical solution of the partial differential equation. On the

other hand, suppose that u0 has a discontinuity at a point x0 ; then, u will have a discontinuity

across the characteristic line x at = x0 issued at x0 . In other words, singularities propagate

along characteristic curves. This very important property is quite general, as it holds for all first

order equations. Obviously, the function defined by (1.3.9) is not a strong solution of (1.3.6) in

O; however, it is a weak solution.

At last, suppose that O is the semi-infinite strip O = {(x, t) : 0 < x < 1, t > 0}. For the sake

of definiteness, assume that a > 0. Then,

(note that the normal vector to O does not exist at the origin, yet this boundary point is an

inflow point for the equation). We prescribe the value u0 = u0 (x) at time t = 0 and the value

g = g(t) at the left endpoint of the interval (0, 1); no condition has to be prescribed at the right

endpoint. Thus, we consider the initial-boundary value problem

ut + aux = 0 0 < x < 1, t > 0 ,

u(x, 0) = u0 (x) 0 < x < 1 , (1.3.10)

u(0, t) = g(t) t>0.

In order to solve this problem, let us fix a point (x, t) O. If x at, then the characteristic passing

through (x, t) meets O at the point (x0 , 0), with x0 = xat; hence, as above, u(x, t) = u0 (xat).

On the other hand, if x < at, then the characteristic passing through (x, t) meets O at the

point (0, t0 ), with t0 = t x/a (see Fig. 1.5); hence, u(x, t) = u(0, t0 ) = g(t0 ) = g(t x/a). We

conclude that the solution of the initial-boundary value problem (1.3.10) is

(

u0 (x at) if x at

u(x, t) =

g(t x/a) if x < at.

Note that if the data u0 and g do not match properly at the origin, a singularity propagates along

the characteristic line x = at.

14 CHAPTER 1. BASIC CONCEPTS

(x, t)

x

t a

0 1 x

If the coefficient a is strictly negative, the boundary condition g is enforced at the right endpoint

of the interval (0, 1).

The concept of characteristic line introduced above is a particular case of the more general

concept of characteristic manifold. A (m 1)-dimensional manifold (a line in two dimensions,

a surface in three dimensions, and so on) contained in O is said non-characteristic for equation

(1.3.1) whenever the following property holds: if one prescribes the value of u on , then u is

uniquely determined by the partial differential equation in a neighborhood of . As a first step,

one aims at determining the gradient of u on ; then, if the manifold, the coefficients and the

data are smoother and smoother, one can differentiate the equation to get derivatives of u on of

higher and higher order; at last, the condition of real analyticity leads to the representation of u

in terms of its Taylor series in a neighborhood of each point in (Cauchy-Kowalewska Theorem).

Confining ourselves to the determination of the gradient of u on , we observe that the directional

derivative of u along any tangential vector to is uniquely determined by the prescribed value of

u therein. Therefore, the differential equation should allow to express the derivative of u along a

non-tangential direction to in terms of the value of u on the manifold. In other words, denoting

by n the normal vector to , one should have a n 6= 0 on . This motivates the following

an=0 on

It is easy to check that characteristic curves lie on characteristic manifolds. Furthermore, the

inflow boundary O of O is obviously a non-characteristic manifold.

The most general linear second order partial differential equation reads as follows:

m

X m

X u

2u

aij + ai + a0 u = f (1.4.1)

xi xj xi

i,j=1 i=1

1.4. LINEAR SECOND ORDER EQUATIONS 15

(the choice of the minus sign in front of the principal part will be motivated in the sequel). We

actually consider the equation in the quasi-divergence form

Xm X m

u

aij + (ai u) + a0 u = f. (1.4.2)

xi xj xi

i,j=1 i=1

As already mentioned above, often the equation is derived in this form; if not, we can transform

(1.4.1) into (1.4.2) by an appropriate modification of lower order coefficients ai (i = 0, 1, . . . , m).

To simplify the notation, let us introduce the square matrix of order m

A := (aij )1i,jm. (1.4.3)

Since uxi xj = uxj xi for any twice continuously differentiable function, it is not restrictive to assume

that aij = aji for all i and j, i.e., to assume that the matrix A is symmetric. Indeed, if the matrix

is not symmetric, we write

1 1

aij uxi xj + aji uxj xi = (aij + aji )uxi xj + (aij + aji )uxj xi ,

2 2

i. e., we replace A by 21 (A + AT ). As before, a = (a1 , . . . , am )T denotes the coefficients of the

first order part. Then, (1.4.2) is compactly written as

Lu = (Au) + (au) + a0 u = f , (1.4.4)

or, equivalently,

Lu = T (Au) + T (au) + a0 u = f . (1.4.5)

A linear second order differential equation can be classified according to the structure of its

principal part. This classification is very important: indeed, the type of the equation influences

the kind of boundary and/or initial conditions which are admissible for the equation, the relevant

properties of the solution, as well as the techniques for solving the equation (analytically or

numerically).

The classification is accomplished by looking at the sign of the eigenvalues of the coefficient

matrix A (recall that A is symmetric, so all its eigenvalues are real). Note that since the coefficients

may depend on x , the type of the equation may vary from point to point. Let us consider A = A(x )

at a fixed point x O.

Three situations are most commonly encountered in applications:

(i) all the eigenvalues of A are not zero, and they all have the same sign; in this case we say

that the operator L (or the equation (1.4.4)) is of elliptic type at x ;

(ii) precisely one eigenvalue of A is zero, while the others have constant sign; in this case we say

that the operator L is of parabolic type at x ;

(iii) all the eigenvalues of A are not zero, and precisely one eigenvalue has a different sign with

respect to the others; in this case we say that the operator L is of hyperbolic type at x .

In two independent variables, this classification is exhaustive (since, by assumption, A cannot

be the null matrix). The terminology comes from the fact that the level curves in the (1 , 2 )-plane

of the associated quadratic form

Q() = T A , = (1 , 2 )T

are ellipses, or degenerate parabolae, or hyperbolae, depending whether the operator L is elliptic,

or parabolic, or hyperbolic.

16 CHAPTER 1. BASIC CONCEPTS

Example 1.4.1. The Poisson equation (1.2.7) is elliptic, the heat equation (1.2.8) is parabolic,

whereas the wave equation (1.2.9) is hyperbolic. Obviously, the type of each equation is the same

at all points in the plane.

Conversely, the Tricomi equation (1.2.10) is of variable type: it is elliptic in the upper half

plane, parabolic on the axis y = 0 and hyperbolic in the lower half plane.

In three or more independent variables, other situations may occur. If A has two or more zero

eigenvalues and the remaing ones are of one sign, we say that the operator is ultra-parabolic. If

two or more eigenvalues are of one sign, whereas two or more remaining ones are of the opposite

sign, we say that L is ultra-hyperbolic. We shall not consider these cases further on.

We now use the classification introduced above to reduce the general second order equation

(1.4.4) to a canonical form. To this end, we shall make the simplifying assumption that the

coefficients of the principal part are constant (otherwise, one can modify the arguments below by

freezing the coefficients in a neighborhood of each point x O).

Denote by i (i = 1, . . . , m) the eigenvalues of A, and let wi be the corresponding eigenvectors,

which form a complete set since A is symmetric. Define the diagonal matrix := diag(1 , . . . , m ),

as well as the orthogonal matrix S := (w1 , . . . , wm ). The eigenvalue-eigenvector relations, written

as AS = S, yield the diagonalization of A

S T A S = . (1.4.6)

Now, let us fix a point x O and let us make the change of independent variable

y = x + S T (x x ).

Denoting by x the gradient in the x -variable and defining y similarly, we have by the chain

rule

x = Sy .

m

X V yj X V m

v

= = sij ,

xi yj xi yj

j=1 j=1

Pm T) x

Pm

since yj = xj + i=1 (S ji i = xj + i=1 sij xi . Substituting into (1.4.5) gives

Lu = Ty S T A S y u + Ty S T (au) + a0 u = f ;

Lu = Ty (y u) + Ty (a u) + a0 u = f. (1.4.7)

m

X 2u

Lu = i + lower order terms.

i=1

yi2

In order to proceed, we consider the three main types of equations introduced above.

1.4. LINEAR SECOND ORDER EQUATIONS 17

If the equation is elliptic, we can assume - possibly after changing the sign of the equation -

that all the eigenvalues of A are positive. Then, we set

1 1

D := diag , . . . ,

1 m

Thus, setting a := Da , (1.4.7) becomes

Lu = z u + Tz (a u) + a0 u = f,

where z is the Laplacian in the z -variable. We conclude that the Laplace operator is the

canonical form of an elliptic operator.

Set n = m 1. Suppose that i > 0 for i = 1, . . . , n, whereas m = 0. If the last component

am of the vector a appearing in (1.4.7) is zero, then the equation does not contain partial

derivatives of u with respect to ym : it is an elliptic equation in n variables, the variable ym

acting only as a parameter. On the other hand, if am 6= 0, we set

1 1 1

D := diag ,..., ,

1 a

n m

and denoting by a the n first components of the vector Da , we transform (1.4.7) into the

form

Dt u z u + Tz (a u) + a0 u = f.

Dt

Set again n = m 1. Suppose now that i > 0 for i = 1, . . . , n, whereas m < 0. Define

1 1 1

D := diag , . . . , ,

1 n m

and perform the same change of variable as in the parabolic case. Setting a := Da , (1.4.7)

becomes

2

Dtt u z u + Tz (a u) + a0 u = f.

We conclude that the wave operator (also termed the DAlembert operator)

2

:= Dtt

18 CHAPTER 1. BASIC CONCEPTS

Let us consider the simple case of a hyperbolic equation, in dimension m = 2. After diagonaliza-

tion, and assuming that the lower order terms are zero, we have

2u 2u

1 2 + 2 2 = f , (1.5.1)

y1 y2

with 1 > 0 and 2 < 0. Let us define a2 := 1 /|2 | > 0; setting x = y1 , t = y2 and g = f /|2 |,

we obtain

2

Dtt u a2 Dxx

2

u = g. (1.5.2)

The equation factorizes as

(Dt + aDx ) (Dt aDx ) u = g, (1.5.3)

which is equivalent to the first order hyperbolic system

(

(Dt + aDx ) w = g (1.5.4)

(Dt aDx ) u = w (1.5.5)

(note that the + and signs can be exchanged in these formulae). Recalling the results of Sect.

1.3, u can be obtained by first integrating (1.5.4) along the characteristics x at = constant,

next integrating (1.5.5) along the characteristics x + at = constant. Actually, the family of lines

x at = constant are called the characteristics of equation (1.5.2). In order to uniquely determine

the solution, one can prescribe a condition on u for each characteristic line at each boundary point

where it enters the region O. Let us detail two examples.

Example 1.5.1. At first, suppose that O is the half-plane {(x, t) : t > 0}. Both characteristics

enter O at each point in O; thus, we prescribe u and a non-tangential derivative of u, such as

the normal derivative ut , therein. Precisely, we consider the initial value problem

utt a2 uxx = 0 x IR , t>0,

u(x, 0) = u0 (x) x IR , (1.5.6)

ut (x, 0) = u1 (x) x IR ,

(where, for simplicity, we have chosen g 0). Taking into account (1.5.4), (1.5.5) and noting

that w(x, 0) = (ut aux )(x, 0) = u1 (x) au0 (x), we first integrate along the characteristics

x at = constant to solve the initial value problem

wt + awx = 0 x IR , t > 0 ,

w(x, 0) = u1 (x) au0 (x) x IR ;

we get w(x, t) = u1 (x at) au0 (x at). Next, we integrate along the characteristics x + at =

constant to solve the initial value problem

ut aux = w x IR, t > 0 ,

u(x, 0) = u0 (x) x IR .

We get Z t

u(x, t) = u0 (x + at) + w(x + at as, s) ds;

0

1.5. BOUNDARY AND INITIAL CONDITIONS. CHARACTERISTICS 19

t t

(x, t)

x + at = x0 x at = x0

x at x + at x x0 x

Figure 1.6: The domain of dependence of a point (x, t) (left) and the domain of influence of a

point (x0 , 0)

substituting the expression of w and making a change of variable in the integral leads to the final

form of the solution:

Z x+at

1 1

u(x, t) = [u0 (x at) + u0 (x + at)] + u1 (s) ds . (1.5.7)

2 2a xat

Rz R0

1

Setting (z) = 21 u0 (z) + 2a 1 1

0 u1 (s) ds and (z) = 2 u0 (z) + 2a z u1 (s) ds, we have

Note that the solution is the superposition of two signals, traveling leftwards and rightwards,

respectively, with speed a and +a; also note that u at (x, t) only depends on the initial data on

the interval [x at, x + at]. If we had considered our equation with a nonzero right-hand side g,

then u(x, t) would have depended on the values of g in the triangle

Indeed, adapting the computations above to the presence of the right-hand side yields

Z x+at

1 1

u(x, t) = [u0 (x at) + u0 (x + at)] + u1 (s) ds

2 2a xat

Z t Z t

+ ds g(x a(2 t s), s) d .

0 s

We call the region T the domain of dependence of the point (x, t) (see Fig. 1.6, left).

Conversely, the initial values at a point (x0 , 0) influence the solution in the angle

A = {(x, t) : x0 at x x0 + at};

this region is called the domain of influence of the point (x0 , 0) (see Fig. 1.6, right).

This simple example shows that a second order hyperbolic equation describes the propagation

and composition of two signals moving at finite speed; the solution depends locally on the data of

the problem (the initial data u0 and u1 , the right-hand side g).

Example 1.5.2. Let us now consider our equation in the semi-infinite strip

20 CHAPTER 1. BASIC CONCEPTS

At each point of the spatial boundary {(0, t) : t > 0} {(1, t) : t > 0}, one characteristics is

entering the domain and one is leaving it, see Fig. 1.7. Thus, one has to prescribe one boundary

condition on u; this can be either the value of u or the value of ux (which is the normal derivative

to O therein). For instance, we can consider the following initial-boundary value problem

0 1 x

utt a2 uxx = 0 x IR , t > 0 ,

u(0, t) = 0 (t) t>0,

ux (1, t) = 1 (t) t>0, (1.5.8)

u(x, 0) = u0 (x) x IR ,

ut (x, 0) = u1 (x) x IR .

In order to motivate the admissibility of the boundary conditions, let us fix a point P0 = (0, t0 )

in O (see Fig. 1.8). If we prescribe u at this point, say u(0, t0 ) = 0 (t0 ), it is convenient to

exchange the signs in (1.5.4), (1.5.5); then, we use the boundary data to integrate u along the

characteristic line x at = at0 entering O at P0 , i.e., we solve

ut + aux = w

u(0, t0 ) = 0 (t0 ),

with w coming from the inside along the characteristic lines x + at = constant.

P0 w

0 x

1.5. BOUNDARY AND INITIAL CONDITIONS. CHARACTERISTICS 21

in (1.5.4), (1.5.5) and we observe that w is known at P0 . Indeed, u has already been determined

for t t0 , so

u(0, t) u(0, t0 )

ut (0, t0 ) = lim ;

tt0 t t0

thus, w(0, t) = ut (0, t0 )a1 (t0 ) and we can integrate w along the characteristic line xat = at0 ,

i.e.,

wt + awx = 0

w(0, t0 ) = ut (0, t0 ) a1 (t0 ).

For an initial-boundary value problem, the domains of dependence and influence are defined

in the obvious way.

as the number of characteristics entering the domain O for increasing t.

It is instructive to consider the situation in which 1 is fixed and 2 tends to 0. In this case, the

speed a tends to infinity, i.e., signals propagate with faster and faster speed. Geometrically, the

slopes of the characteristic lines in the (x, t)-plane tend to 0, and the domains of dependence and

influence of any point get wider and wider. In the limit, eq. (1.5.1) becomes the elliptic equation

2u

1 =f

x2

in the sole space variable y1 = x. The solution at each point x depends on the values of the data

f at all points x in the domain, as well as on the boundary data at all boundary points.

If a low order term is present in the equation, i.e., if the equation is

2u 2u u u

1 2 + 2 2 + a1 + a2 =f

x t x t

with a2 6= 0, the limit equation for 2 0 is parabolic; the solution at each point (x, t) in the

domain depends on all the values of the data f and the boundary data for all t < t, as well as on

the initial condition u0 at t = 0. Propagation of signals takes place with infinite speed.

At last, we briefly deal with the concept of characteristic manifold. For a second order partial

differential equation, a manifold O is non-characteristic if the prescription of u and u on

uniquely determines the Hessian of u (i.e., the set of all its second order partial derivatives)

therein, via the differential equation.

Suppose that the manifold is described by an implicit equation (x ) = 0, for a smooth . Fix

a point x on and let

(x )

n=

k(x )k

be the normal vector to at x . Let us make a change of independent variable y = x + RT (x x ),

1

such that the last coordinate direction is along n. Setting (y ) = x + RT (y x ) , we

have R1 x = y , so we choose R such that R1 n = em = (0, . . . , 0, 1)T . The differential

equation in the new coordinates becomes

Ty RT A R y u + lower order terms = f.

22 CHAPTER 1. BASIC CONCEPTS

We note that all second order derivatives of u except uym ym are determined at x by the values

of u and u on . Thus, in order to get the value of uym ym , we must have

RT A R mm = eT T T

m R A R em = n A n 6= 0.

satisfies

nT A n = 0 on

is called a characteristic manifold for equation (1.4.1).

For instance, the characteristic manifolds of the wave equation (1.5.2) are defined by the

equation a2 n2x n2t = 0 (where n = (nx , nt )T ), i.e., they are precisely the lines x at = constant.

The characteristic manifolds for the heat equation (1.2.8) satisfy n2x = 0, i.e., they are the lines

t = constant.

Finally, any elliptic equations has no (real) characteristic manifold. This means that, under

appropriate regularity conditions, the Cauchy problem

Lu = f in O ,

u = u0 on ,

u

= u1 on ,

n

is always uniquely solvable in a neighborhood O of any (m 1)-dimensional manifold . However,

the Cauchy problem is not well-posed for an elliptic equation. This means that arbitrarily small

changes in the data u0 and u1 may lead to arbitrarily large changes in the solution u, as the

following example shows.

sin nx u u

u(x, 0) = , (x, 0) = (x, 0) = 0 ,

n n y

for a fixed n > 0 (thus, u(x, y) = un (x, y)). The solution can be found by the ansatz

sin nx

u(x, y) = u(y),

n

which reduces the problem to a second order ordinary differential equation with two initial condi-

tions:

1

u nu = 0

n

u(0) = 1

u (0) = 0.

The result is

eny + eny

u(y) = = cosh ny,

2

1.6. EXERCISES 23

sin nx

u(x, y) = cosh ny.

n

As n , the initial data converge to 0 uniformly in n, whereas u becomes arbitrarily large in

an arbitrarily small neighborhood of O.

For this reason, an elliptic equation is more appropriately supplemented by one boundary

u

condition, involving u and/or n , at each point of the boundary O of the region O where the

equation is set. In this way, one obtains a well-posed problem, as we shall see in Chapter 4.

1.6 Exercises

1.1. Consider the transport equation

u u

+2 =0

t x

in the half-plane {(x, t) : t > 0}, with the initial condition

(

3 if x < 0

u(x, 0) = u0 (x) =

1 if x > 0.

3 if t > 21 x

u(x, t) =

1 if t < 12 x

is a weak, not classical, solution of the equation.

u u

+x =1

t x

and:

(ii) solve the initial value problem in O = IR (0, +) with the condition u(x, 0) = u0 (x);

(iii) solve the initial-boundary value problem first for x [0, 1] and then for x [1, 2] with the

further condition u = g(t) on the inflow boundary.

u u

+u =0

t x

is the simplest example of nonlinear transport equation.

(ii) Deduce that the characteristics are straight lines in the half plane {(x, t) : t > 0}.

(iii) Suppose the initial datum u(x, 0) = u0 (x) is prescribed for every x IR; find the slope of

the characteristics.

24 CHAPTER 1. BASIC CONCEPTS

2u 2u 2u 2u

(i) + 3 + + 4 =f

x2 xy yx y 2

2u 2u

(ii) + y = g.

x2 xy

Chapter 2

Theory of Distributions

The theory of distributions was created by Laurent Schwartz in 1944; its main purpose is to extend

the results which hold for integrable and differentiable functions to those functions that do not

satisfy the necessary conditions of classical regularity.

Let O be an open set in IRm ; we recall that the set D(O) has been previously defined as

where supp denotes the support of , i.e., the closure of the set of all points x in O such that

does not vanish on them:

supp = {x O | (x ) 6= 0};

it is easy to verify that D(O) is a linear space.

Let us now introduce the following notion of convergence in D(O):

(i) there exists a compact set K O which contains all the supports of n and ;

(ii) for all multi-integers INm , the sequence {D n }n0 converges to D uniformly on K,

i.e.,

n

kD n D k,K 0.

T : D(O) IR

such that if {n }n0 converges to in D(O) then {T (n )}n0 also converges to T () in IR when

n .

The set of all distributions on O is a linear space denoted by D (O). Moreover, the notation

hT, i is often used instead of T () and it is called a duality form.

25

26 CHAPTER 2. THEORY OF DISTRIBUTIONS

Example 2.1.3. Let f be a real-valued and Riemann (or Lebesgue)-integrable function on O; let

us set Z

hTf , i := f (x )(x ) dx D(O)

O

and let us verify that Tf is a distribution. To do this, we have to check the properties of the

previous definition; in particular:

(i) Tf is certainly a linear form because it is real-valued and the integral is a linear operator;

(ii) suppose {n }n0 D(O) is a sequence such that kn k,K 0 when n for a

certain D(O); then

Z

hTf , n i hTf , i = hTf , n i = f (x )[n (x ) (x )] dx

O

and so

Z

|hTf , n i hTf , i| |f (x )| |n (x ) (x )| dx

K Z

kn k,K |f (x )| dx =

K

n

= Gkn k,K 0

R

where G = K |f (x )| dx is a finite constant that comes from the hypothesis that f is inte-

grable on K. Thus we have hTf , n i hTf , i when n .

Z Z

f (x )[n (x ) (x )] dx = f (x )[n (x ) (x )] dx

O K

because the supports of all the n s and of are contained in K, so the integral vanishes on O\K.

This imply that only a local integrability of f on subdomains of O, and not on the whole set O,

is needed to define the distribution Tf .

Throughout this chapter, we shall refer to this type of distribution as a function-like distri-

bution.

Example 2.1.4 (The Dirac delta). Consider a point x 0 O; we introduce now the following

form

hx 0 , i := (x 0 ) D(O)

and we want to verify that it is a distribution in the sense of Definition 2.1.2.

(ii) Let us suppose that n in D(O); by Definition 2.1.1, for all we have a uniform

convergence of D n to D , then for || = 0 it follows

n

max |n (x ) (x )| 0

x K

and so, if x 0 K:

n

|hx 0 , n i hx 0 , i| = |n (x 0 ) (x 0 )| max |n (x ) (x )| 0.

x K

2.1. BASIC DEFINITIONS 27

1 1

2n 2n x

Such a distribution is called the Dirac delta on the point x 0 ; it is possible to show (see Exercise

2.1) that it is not a function-like distribution, i.e., it does not exist any function f such that the

action of x 0 on a test function D(O) can be expressed as the integral on O of f versus .

After introducing the notion of convergence in D(O), it would be useful to provide a similar

tool for the space D (O) too. This is accomplished by the following

Definition 2.1.5. Let T, Tn D (O), n 0; the sequence {Tn }n0 is said to converge to T in

the sense of D (O) if

n

hTn , i hT, i

for every D(O).

This definition leads us to an important characterization of the Dirac delta. Let us set O = IR

and T = 0 , hT, i = (0) for all D(IR); then, for every n > 0, let us define the function (see

Figure 2.1) (

1

n if |x| 2n

fn (x) =

0 otherwise.

Z Z 1

2n

fn (x) dx = n dx = 1

1

IR 2n

does not depend on n: every function fn has therefore the same unitary area on IR. If we now

consider the family of distributions Tfn , we have:

Z Z 1

2n 1

hTfn , i = fn (x)(x) dx = n (x) dx = n (xn ) = (xn )

IR 1

2n n

1 1

where xn is a point in the interval 2n , 2n whose existence is guaranteed by the Integral Mean

Theorem. It is clear that xn 0 when n ; then using the continuity of gives

n

hTfn , i (0) = h0 , i.

28 CHAPTER 2. THEORY OF DISTRIBUTIONS

Since this argument holds for every D(IR), we conclude that Tfn 0 in the sense of D (IR).

This show that, although the Dirac delta cannot be represented by a classical function, it can

nevertheless be obtained as a limit of classical functions in the sense of Definition 2.1.5.

R

In general, it is easy to check that any sequence {fn } of integrable functions satisfying R fn (x) dx =

1 and supp fn B(0, rn ) with rn 0 as n , converges to 0 in D (IR) as n .

Definition 2.1.6. A distribution T is said to be of finite order if there exist r IN and a constant

Cr > 0 such that

D(O), |hT, i| Cr max kD k,O .

||r

The smallest r for which this condition holds is called the order of the distribution.

Z Z

|hTf , i| = f (x )(x ) dx kk,O |f (x )| dx

O O

Z

0 |f (x )| dx = C < +

O

so

|hTf , i| Ckk,O .

It is possible to verify that this is also the order of the Dirac delta x 0 .

Definition 2.1.8. Let T D (O); the support of T is the smallest closed set K O such that

This definition states that the support of a distribution T is strictly related to those of test

functions. More in detail, the support K of T is the smallest closed set in O that has the following

property: every test function that vanishes on the whole K, i.e., such that its support does not

intersect K, sees T as zero.

For instance, if we take T = x 0 , x 0 O, we find supp x 0 = {x 0 } because every test function

whose support does not contain x 0 is such that (x 0 ) = 0 and so hx 0 , i = 0.

As another example, let us consider an integrable function f with a compact support in O;

then supp Tf = supp f .

Example 2.1.9. Consider an open set O in IRm and let be a closed (m 1)-dimensional regular

manifold contained in O; let g be an integrable function defined on . Then, the distribution ,g

defined as Z

h,g , i = g()() d D(O),

2.2. DERIVATIVES OF DISTRIBUTIONS 29

In this section, the main results from the differential theory of distributions are exposed. In

particular, we shall see, with the aid of many examples, in which sense such a theory represents

a generalization of the classical one and what meaning has to be given to the word derivative

referred to a distribution.

Let us start with this basic definition.

Definition 2.2.1. Let INm and T D (O); the partial derivative of T of order is the

distribution D T whose action on a test function D(O) is defined as

hD T, i = (1)|| hT, D i.

We can immediately observe that, in the sense of this definition, all the distributions are

infinitely differentiable, since the derivative is moved on the test function which is of class C (O).

The following example will explain the reason of such a definition and where it comes from.

Example 2.2.2. Let f C1 (O) and consider the distribution T = Tf ; in order to calculate its

derivative Di Tf , we set = (0, . . . , 0, 1, 0, . . . , 0) (where the only component of the multi-integer

different from zero is the i-th) and then we apply the Definition 2.2.1:

Z

hDi Tf , i = hTf , Di i = f (x ) (x ) dx =

O xi

Z

f

= (x )(x ) dx = hTDi f , i

O xi

for all D(O); we recall that in applying the integration-by-parts formula no boundary term

appears, since a test function vanishes in a neighborhood of O.

We conclude that Di Tf = TDi f , i.e., the partial derivative with respect to xi of the distribu-

tion based on the function f is the distribution based on the function Di f , which exists in the

classical sense under the hypothesis f C1 (O). As we have just seen, this result follows from

the integration-by-parts formula and it allows us to calculate the derivatives of a function-like

distribution in a somewhat classical way.

(

2x if x 0

f (x) =

x if x < 0

which is not differentiable in the classical sense because of the singularity at the origin. Neverthe-

less, in the distributional sense we have:

Z

h(Tf ) , i = hTf , i = f (x) (x) dx =

IR

Z 0 Z +

= x (x) dx 2x (x) dx =

0

Z 0 Z +

= (x) dx + 2(x) dx =

0

Z

= g(x)(x) dx = hTg , i D(IR)

IR

30 CHAPTER 2. THEORY OF DISTRIBUTIONS

2 if x > 0

g(x) =

1 if x < 0;

then (Tf ) = Tg or, as often one writes, f = g in the sense of distributions.

Note that the derivative of Tf is itself a function-like distribution; this depends strictly on the

fact that f is continuous on IR.

Example 2.2.4. Let us now consider the function

(

2x + 1 if x > 0

f (x) =

x if x < 0

which is discontinuous at the origin. In this case we have:

Z

h(Tf ) , i = hTf , i = f (x) (x) dx =

IR

Z 0 Z +

= x (x) dx (2x + 1) (x) dx =

0

Z 0 Z + +

= (x) dx + 2(x) dx (2x + 1)(x) =

Z 0 0

IR

where g is defined as in the example above. Then

h(Tf ) , i = hTg , i + h0 , i D(IR)

and consequently (Tf ) = Tg +0 , which is no longer a function-like distribution although Tf is.

In general, if one has (

f+ (x) if x > x0

f (x) =

f (x) if x < x0

with f+ C1 [x0 , +), f C1 (, x0 ], then

(Tf ) = Tg + |[f ]|x=x0 x0 , (2.2.1)

where (

f+ (x) if x > x0

g(x) =

f (x) if x < x0

and |[f ]|x=x0 denotes the jump of f at the point x0 . Therefore, (Tf ) is a function-like distribution

if, and only if, f is continuous at x0 ; in fact, in this case |[f ]|x=x0 = 0, which eliminates the delta

from the expression (2.2.1).

Example 2.2.5. The Heaviside function is defined as

(

1 if x > 0

H(x) =

0 if x < 0;

from (2.2.1) it follows (TH ) = 0 . The Heaviside function is then a primitive of the Dirac delta

in the sense of distributions.

More often one writes H = 0 , where the derivative is, of course, intended in the distributional

sense.

2.3. STUDY OF THE LAPLACE OPERATOR IN D (O) 31

Example 2.2.6. Consider the Dirac delta 0 D (IR) and let D(IR); from Definition 2.2.1

one has

h0 , i = h0 , i = (0)

h0 , i = h0 , i = (0)

..

.

(k)

h0 , i = = (1)k (k) (0), k IN.

Example 2.2.7. The multidimensional counterpart of the general situation considered in Example

2.2.4 is as follows. Let O an open set in IRm and let be an (m 1)-dimensional regular manifold

contained in O, which splits O as O O+ , with O open disjoint sets such that O O+ = .

Let the function f satisfy (

f (x ) if x O

f (x ) =

f+ (x ) if x O+

with f+ C1 (O+ ), f C1 (O ). Then, for any i = 1, . . . , m, one has

Di (Tf ) = Tgi + ,hi (2.2.2)

where (

Di f (x ) if x O

gi (x ) =

Di f+ (x ) if x O+

and

hi (s) = |[f ]|s ni (s), s ,

where |[f ]|s denotes the jump of f at the point s in going from O to O+ , and ni is the i-th

component of the normal unit vector to pointing from O to O+ .

We prove the result in the particular case in which f is the Heaviside function associated with

the given partition of O, i.e., (

0 if x O

H(x ) =

1 if x O+

Let us compute Di (TH ) in the sense of distributions. Using the divergence theorem (see (3.1.3)),

we have

Z Z Z

hDi (TH ), i = H(x ) (x ) dx = (x ) dx = ni d = h,ni , i

O xi O+ xi

We refer to Exercise 2.6 for the proof of the general result.

In this section we study the Laplacian

Xm

2

=

i=1

x2i

as an operator into the space of distributions D (O); in particular, we are interested in those

functions g : O IRm IR whose Laplacian is the Dirac delta 0 on the origin.

32 CHAPTER 2. THEORY OF DISTRIBUTIONS

g = 0 in D (O)

Let us start with m = 1 (dimension 1); if we take the function u(x) = |x|, it is easy to verify

that u (x) = sign(x) and thus u (x) = 20 , as it immediately follows from (2.2.1). Hence, the

function g(x) = 21 u(x) = 12 |x| is a fundamental solution of the Laplacian on IR.

Let us now consider m = 2; in this case, it is convenient to use the polar coordinates defined

by the transformation

(r, ) 7 (x, y) = (r cos , r sin );

since we can think of every function u = u(x, y) as u(x, y) = u((r, )) = U (r, ), we have the

relationship u(x, y) = U (r, ) and consequently (x, y) u = (r, ) U , where (x, y) and (r, ) denote

the Laplacian in cartesian and polar coordinates respectively, with (see Exercise 2.8)

2 1 1 2

(r, ) =+ + . (2.3.1)

r 2 r r r 2 2

p

If we take the function u(x, y) = log x2 + y 2 = log r and we set (x, y) = 6 (0, 0), we obtain

from (2.3.1)

1 1

u = (r, ) log r = 2 + 2 = 0;

r r

hence, log r is a harmonic function in the classical sense everywhere in the plane except at the

origin.

Let us now calculate u in the sense of distributions; taking D(IR2 ) we have

hu, i = hu, i =

Z Z

= log r dx dy = lim log r dx dy

IR2 0+ IR2 \B(0, )

where B(0, ) is the open ball of radius > 0 centered at the origin. Applying the integration-by-

parts formula gives

Z Z

hu, i = lim log r dx dy + log r d =

0+ n

Z r> Z r=

Z

= lim log r dx dy log r d + log r d =

0+ r> r= n r= n

Z Z

= lim log r d + log r d

0+ r= n r= n

where the result log r = 0 out of the origin has been used.

Since B(0, ) is a circle, the normal vector of its circumference r = is a radial vector, which

allows us to write

d 1

log r = log r =

n dr r

2.3. STUDY OF THE LAPLACE OPERATOR IN D (O) 33

where the minus sign depends only on the fact that n log r = log rn should be negative because

the two vectors log r and n point in opposite directions. Thus

Z Z Z

1 1

log r d = d = 2 d.

r= n r= 2 r=

1

R

Note that 2 r= d is the mean value that takes along the circumference r = ; since

is continuous, it follows

Z

1

lim 2 d = 2(0, 0).

0+ 2 r=

Moreover Z Z

log r d = log d

r= n r= n

and it results

Z Z Z

d | n| d

n n d =

r=

Zr= r=

Z

kk d max kk d = 2 max kk;

r= (x, y)IR2 r= (x, y)IR2

in the third passage, the Cauchy-Schwartz inequality has been used within the fact that knk = 1

(here k k denotes the Euclidean norm in IR2 ). Since

2 max kk = M

(x, y)IR2

Z

0+

log M | log | 0

r= n

and finally

hu, i = 2(0, 0) = 2h0 , i D(IR2 )

that is

u = 20 in D (IR2 ).

1 1

p

Hence, the function g(x, y) = 2 u(x, y) = 2 log x2 + y 2 is a fundamental solution for the

Laplacian on IR2 .

In three dimensions, with the aid of the spherical coordinates, it can be found that the function

1

u(x, y, z) = p

x2 + y2 + z2

is such that u = 40 ; it is therefore proportional to a fundamental solution on IR3 .

In general, we have the following expressions for the fundamental solutions of the Laplacian:

1

r m=1

2

1 v

log r m=2 um

2 uX

g(x ) = r = kx k = t x2i (2.3.2)

1 1

m = 3 i=1

4 r

1 1

m2 m 4

(m 2)m r

34 CHAPTER 2. THEORY OF DISTRIBUTIONS

2 m/2

where m = is the surface area of the unit sphere in IRm .

(m/2)

It is obvious that adding any harmonic function to g, i.e., a function v such that v 0,

leads to another fundamental solution of the Laplacian. Actually, we are more interested in the

existence rather than in the uniqueness of the fundamental solutions, since their importance is

due to the fact that they provide a powerful tool for solving the following more general matter:

find u such that u = f in O, where f is a given bounded function with a compact support and

integrable on O (i.e. f L1 (O)).

Note that, given any function g such that g = 0 , the new function

has the following property: if we denote by x the Laplacian with respect to the variable x , then

x G = y in D (IRm ),

because the singularity of g has now been moved from the origin to the point y .

Let us set

Z

u(x ) := (f g)(x ) = g(x y )f (y) dy =

Z O

= G(x , y )f (y ) dy;

O

Z

hu, i = hu, i = u(x )(x ) dx =

Z Z O

= G(x , y )f (y )(x ) dx dy =

O O

Z Z

= f (y) G(x , y )(x ) dx dy ;

O O

Z

G(x , y )(x ) dx = hG, i = hx G, i = hy , i = (y )

O

and then Z

hu, i = f (y)(y ) dy = hf, i;

O

since this argument holds for every test function D(O), we conclude that such a u is a solution

of the elliptic equation u = f in the sense of distributions.

2.4 Exercises

2.1. Prove that the Dirac delta is not a function-like distribution, i.e., that it does not exist any

integrable function f : O IRm IR such that

Z

hx0 , i = f (x )(x ) dx , D(O).

O

2.4. EXERCISES 35

2.2. Consider an open set O in IRm and let be a (m 1)-dimensional regular manifold contained

in O. Moreover, let g be an integrable function defined on ; prove that the formula

Z

hT,g , i = g() () d

n

2

0 if x < 0 or x > n

fn (x) = n2 x if 0 x n1

2n n2 x if n1 < x n2 ;

n

2.5. Let be the straight line in the plane having equation y = 2x. Define then the distribution

D (IR2 ) such that Z

h , i = () d

2

for every test function D(IR ).

u

= in D (IR2 ).

x

v v

+ = in D (IR2 ).

x y

u

= 0 in D (IR2 ).

x

Which is the support of u?

2.8. Prove that the Laplacian in polar coordinates is given by equation (2.3.1).

36 CHAPTER 2. THEORY OF DISTRIBUTIONS

Chapter 3

Sobolev Spaces

3.1 Motivation

In order to motivate the introduction of the Sobolev space H1 (), let us consider the following

Dirichlet boundary-value problem for a general second-order elliptic operator Lu:

(

Lu = (Au) + (au) + a0 u = f in

(3.1.1)

u = 0 on .

Here, A, a, a0 and f are known functions defined in ; precisely, A takes its values in the space

of symmetric and positive-definite matrices of order n, a is a vector-valued function, whereas a0

and f are scalar functions.

We aim at giving a weak (or integral, or variational) formulation of this problem, which

corresponds to the general form (1.2.16). At the beginning, we will proceed in a formal manner,

assuming that all mathematical operations are permitted; then, step by step, we will envisage a

set of assumptions on the data of the problem (the coefficients of the operator, the right-hand

side, the domain) which make the resulting formulation mathematically rigorous.

The starting point consists of multiplying the first equation in (3.1.1) by a test function v and

integrating over , to get

Z Z Z Z

(Au)v + (au)v + a0 uv = fv . (3.1.2)

Next, we perform an integration-by-parts in the first and second term on the left-hand side.

Precisely, we invoke the divergence theorem

Z Z

F= Fn, (3.1.3)

where F is a vector field and n is the unit vector which is normal to and pointing outwards,

as well as the differentiation rule for a product

(v) = ( ) v + v , (3.1.4)

where is a vector field and v is a scalar function. Applying (3.1.3) and (3.1.4) to F = v with

= Au, we obtain

Z Z Z

(Au)v + (Au) v = n (Au) v ;

37

38 CHAPTER 3. SOBOLEV SPACES

u

= n (Au) (3.1.5)

nA

u

(which coincides with the normal derivative = n u when A is the identity matrix), we get

n

Z Z Z

u

(Au)v = (Au) v v. (3.1.6)

nA

Z Z Z

(au)v = u a v + a n uv . (3.1.7)

Z Z Z Z Z Z

u

(Au) v u a v + a0 uv v+ a n uv = fv . (3.1.8)

nA

Now, we observe that u is required to vanish on ; therefore, from now on, we will require

that our test functions v vanish on , too (note that functions in D() do satisfy this condition).

Then, (3.1.8) simplifies as

Z Z Z Z

(Au) v u a v + a0 uv = fv . (3.1.9)

Note that this equation only involves first-order partial derivatives of u and v.

Next, we make assumptions on the functions appearing in (3.1.9), so that all integrals therein

are guaranteed to be meaningful and finite. On the left-hand side, we have integrals of products

of three functions, such as

Z Z Z

u v v

aij or ai u or a0 uv ,

xj xi xi

whereas the right-hand side is the integral of the product of two functions. Thus, we set ourselves in

the framework of the Lebesgue Integration Theory, which, in particular, ensures that the product

of two functions is integrable in , i.e., L1 () if Lp () and Lp () with

p, p [1, ] satisfying p1 + p1 = 1; furthermore, the following Holder inequality holds:

Z Z Z 1/p Z 1/p

|| ||p

||p

= kkLp () kkLp () (3.1.10)

R 1/p

(if p = , the term ||p has to be replaced by ess sup ||, and similarly if p = ). This

result extends to the product of three functions, i.e., L1 () if Lp (), Lp () and

Lp () with p, p , p [1, ] satisfying p1 + p1 + p1 = 1; in this case, one has

Z

kkLp () kk p kk p . (3.1.11)

L () L ()

The structure of the integrals in (3.1.9) suggests to work in a Hilbertian setting, i.e., to assume that

u, v and their first derivatives belong to L2 (). More precisely, the previous results tell us that

3.2. THE SPACE H1 () 39

R R

f v is well-defined if f and v L2 (); a0 uv is well-defined if a0 L () and u, v L2 ();

R v

an integral of the form ai x u is well-defined if ai L (), u L2 () and x

v

L2 (); finally,

R i

u v

i

an integral of the form aij x j xi

is well-defined if aij L (), xu

j

and x v

i

L2 (). In

conclusion, if we assume that

then u and v should belong to L2 () together with all their first-order partial derivatives. Such

derivatives have to be considered in the sense of distributions, since u and v are merely L2 -

integrable functions, and not classical differentiable functions.

This leads us to introduce the Sobolev space H1 () and, subsequently, its closed subspace

1

H0 () of the functions vanishing on . This will be the appropriate space for setting the weak

formulation of problem (3.1.1) and for studying its well-posedness.

Motivated by the previous discussion, we introduce the space H1 () as follows.

Definition 3.2.1. H1 () is the subspace of L2 () of the functions whose first-order partial deriva-

tives, in the distributional sense, belong to L2 (), i.e.,

v

H1 () = {v L2 () : L2 () for 1 i n} = {v L2 () : v (L2 ())n } .

xi

n

X

u v

(u, v)H1 () = (u, v)L2 () + , = (u, v)L2 () + (u, v)(L2 ())n ,

xi xi L2 ()

i=1

n !1/2

X v 2 1/2

kvkH1 () = kvk2L2 () + = kvk2L2 () + kvk2(L2 ())n .

xi 2

i=1 L ()

We point out that the requirement v/xi L2 () means that there exists gi L2 () such that

Tv /xi = Tgi in the sense of distributions, i.e.,

Z Z

h Tv , i = hTv , i= v = gi = hTgi , i D() ;

xi xi xi

By the very definition of the norm in H1 (), one has kvkL2 () kvkH1 () for all v H1 (),

i.e., the inclusion H1 () L2 () is continuous.

Next property is one of the fundamental properties of H1 ().

Proof. Let {vk }k0 be a Cauchy sequence in H1 ()-norm, i.e., > 0, k IN such that

, m > k one has kv vm kH1 () < . This immediately implies that each sequence {vk }k0 ,

40 CHAPTER 3. SOBOLEV SPACES

{vk /xi }k0 for i = 1, . . . , n, is a Cauchy sequence in L2 (). By the completeness of this space,

there exist functions v and gi , i = 1, . . . , n, belonging to L2 (), such that

vk

lim vk = v , lim = gi , i = 1, . . . , n ,

k k xi

in L2 (). The property is proven if we prove that v/xi = gi for i = 1, . . . , n. This follows from

Z Z Z

v

h , i = v = lim vk = lim vk

xi xi k xi k xi

Z Z Z

vk vk

= lim = lim = gi = hgi , i D() .

k xi k xi

Next property, which we state without proof, is important both from the theoretical and

the constructive/numerical point of view; indeed, it guarantees that functions in H1 () can be

approximated arbitrarily well by functions belonging to a sequence of finite dimensional subspaces.

Property 3.2.3. H1 () is separable, i.e., it contains a sequence {vk }k0 which is dense in it.

Let us improve our knowledge of the space H1 () by observing that it contains classical dif-

ferentiable functions. Indeed, if is bounded, any function v C 1 () belongs to H1 (), and one

has !1/2

n

X

v
2

kvkH1 () ||1/2 2

kvkC 0 () +
(n + 1)||1/2 kvkC 1 () .

xi
0

i=1 C ()

belongs to H1 (). In particular, for any open set , one has:

So far, we have seen that sufficiently smooth functions, in a classical sense, belong to H1 ().

On the other hand, H1 () also contains piecewise smooth functions, provided they are globally

continuous. The following result illustrates the situation.

Property 3.2.6. Let be a bounded open set, which is divided into two open subsets and +

by a smooth (n 1)-dimensional manifold . Given two functions v C 1 ( ), the function v

defined as (

v (x) if x ,

v(x) =

v+ (x) if x + ,

belongs to H1 () if and only if v is continuous across .

v

(x ) if x ,

xi

gi (x ) =

v

+ (x ) if x + ,

xi

3.2. THE SPACE H1 () 41

Tv = Tgi + ,[v] ni ,

xi

where [v] is the jump of v across and ni is the i-th component of the normal vector n to . The

result follows from the observation that gi L2 (), whereas ,[v] ni 6 L2 () unless [v] ni 0.

The previous result has a strong practical impact, as it guarantees that one can use continu-

ous, piecewise polynomial functions in order to approximate the solution of second-order elliptic

problem; the finite element method relies precisely on this property.

One may wander if a function belonging to H1 () is more regular than just an L2 ()-function,

for instance if it is continuous. First, let us clarify the real meaning of the statement a function

v H1 () is continuous. Indeed, H1 () is a subspace of L2 (), which according to the Lebesgue

integration theory rigorously speaking does not contain functions but classes of equivalence of

functions, two functions in the same class differing only on a zero-measure subset of . Then, the

statement above means that in the equivalence class of v there exists a function, say v, which is

continuous. For simplicity, in the sequel of this book we will not distinguish a function and the

equivalence class which contains it. The following result gives a positive answer in dimension 1.

Property 3.2.7. If = I is a bounded interval, then H1 (I) C0,1/2 (I) with continuous injection,

where C0,1/2 (I) is the space of the Holder continuous functions of exponent 1/2 in I.

Proof. Let us fix v H1 (I) and let us set g = v L2 (I). Let us define the function

Z x

w(x) = g(s) ds ,

x0

where x0 is any fixed point in I. Since v = w in the distribution sense, there exists a constant C

such that v(x) = w(x) + C in I. Thus, for any two points x1 , x2 I,

Z x1

v(x1 ) v(x2 ) = w(x1 ) w(x2 ) = g(s) ds ,

x2

Z Z 1/2 Z 1/2

x1 x1 x1

|v(x1 ) v(x2 )| = 1 g(s) ds 1 ds

2

g (s) ds |x1 x2 |1/2 kgkL2 (I) .

2

x2 x2 x2

This precisely means that v is Holder continuous of exponent 1/2 in I, and that

|v(x1 ) v(x2 )|

|v|C0,1/2 (I) := sup kv kL2 (I) . (3.2.1)

x1 ,x2 I |x1 x2 |1/2

Z Z

v(y) dy = 1 v(y) dy |I|1/2 kvk 2 . (3.2.2)

L (I)

I I

Z

1

v(x) = v(y) dy

|I| I

42 CHAPTER 3. SOBOLEV SPACES

v(x) v(x)

v(x) = v(x) + |x x|1/2 ,

|x x|1/2

and we apply (3.2.1) and (3.2.2) to get

|v(x) v(x)| 1/2

|v(x)| |v(x)| + |I| |I|1/2 kvkL2 (I) + |I|1/2 kv kL2 (I) ,

|x x|1/2

which easily implies kvkC0 (I) C1 kvkH1 (I) . In conclusion, this estimate and (3.2.1) yield

counterexample shows. Consider the disc = {(x, y) IR2 : r = x2 + y 2 < 1/2} and the

function v(x, y) = | log r| for some > 0. Obviously, v is unbounded as r 0, hence, in

particular, it is not continuous at the origin. Let us show that v H1 () if and only if < 1/2.

We have Z Z Z 2 1/2

v 2 dxdy = | log r|2 r drd < + for any > 0 .

0 0

On the other hand,

v x v y

= | log r|1 2 , = | log r|1 2 ,

x r y r

whence

Z 2 2 Z Z 1/2

v v

+ dxdy = 2 | log r|22 1 dxdy = 22 1

| log r|22 dr < + iff < 1/2 .

x y r 2 r

0

We will see later on (Thm. 3.8.2) that functions in H1 (), although not necessarily continuous,

belong to some space Lp () with p > 2 depending on n.

Another fundamental result is the following one. Let D() be the space of the C -functions

defined in , whose support is compact and contained in (thus, they are allowed to be nonzero

on ). Note that D() = C () if is bounded, whereas D() = D() iff = IRn . The space

D() can equivalently be defined as the space of the restrictions to of the functions in D(IRn ).

Property 3.2.8. D() is a dense subspace of H1 ().

We will give some ideas of the proof, in some particular cases, later on. The property states

that any function in H1 () can be approximated arbitrarily well by smooth classical functions.

This will allow us to pass to the limit and extend results which are well-known for classical

functions to analogous results for functions in H1 ().

The definition of H1 () can be generalized, by considering for any m 2 the set of all functions

v L2 () such that all their distributional partial derivatives D v of order || m belong to

L2 (). Thus, we set

Hm () = {v L2 () : D v L2 () for || m}

3.4. THE SPACES HS (IRN ) 43

X

(u, v)Hm () = (D u, D v)L2 () ,

||m

X

kvkHm () = kD vk2L2 ()

||m

(we use here the convention that D0v = v). In this way, we obtain a separable Hilbert space which

enjoys properties similar or equal to those seen for H1 (): for instance, it contains D() as a dense

subspace and, if is bounded, it contains Cm () but also those functions of Cm1 () which are

piecewise Cm -differentiable. Furthermore, all functions in Hm () enjoy classical differentiability

of some order < m (which depends on the space dimension n); for instance, in dimension n = 2,

the space H2 () is contained in C0 () with continuous inclusion (but not in C1 ()). The precise

result will be given in Thm. 3.8.2.

We thus have a scale of function spaces, in which smoothness is measured in a weak, integral

sense; each space is strictly contained in all the spaces of lower index, with continuous inclusion.

Such a scale is the counterpart of the classical scale of spaces Cm (), in which smoothness is

measured in a strong, pointwise sense. Precisely, the two sequences of spaces satisfy

Hm+1 () Hm () Hm1 () H1 () H0 () = L2 ()

Cm+1 () Cm () Cm1 () C1 () C0 ()

and if is bounded each space of the lower sequence is contained with continuous inclusion in

the space above it in the upper sequence. Working in the Sobolev scale rather than in the classical

one is more appropriate for handling the weak, or integral, formulation of an elliptic boundary

value problem; in particular, the Sobolev scale consists of Hilbert spaces, whereas the classical

scale consists merely of non-reflexive Banach spaces.

A further generalization comes from replacing L2 () by some Lp () with p [1, +] in the

definition of the Sobolev space. Thus, we set

Wm,p () = {v Lp () : D v Lp () for || m}

equipped with the norm

1/p

X

kvkWm,p () = kD vkpLp () .

||m

Such a space is a Banach space, which as Lp () is reflexive if 1 < p < + and is non-reflexive

if p = 1 or p = +. Note that Wm,2 () = Hm (). Sobolev spaces of summability index p 6= 2

play a crucial role in studying nonlinear partial differential equations.

The study of Sobolev spaces is particularly interesting and important when the domain is the

full space IRn . In this section, we first provide a characterization of H1 (IRn ) by means of the

Fourier transform. Next, we consider the spaces Hm (IRn ) for m > 1 and we extend the definition

of Sobolev spaces to the case where the index is any real number. Finally, we sketch the proof of

Property 3.2.8 in the present situation, i.e., when the boundary is empty.

44 CHAPTER 3. SOBOLEV SPACES

We recall that the (continuous) Fourier transform

Z

1

v(x ) 7 v() = v(x ) eix dx

(2)n/2 IR n

Z

1

v() 7 v(x ) = v() e ix d ;

(2)n/2 IRn

more precisely, the transform is an isometry, i.e., for all v, w L2 (IRn ) one has

Z Z

v(x )w(x ) dx = v()w() d , whence, kvkL2 (IRn ) = kvkL2 (IRn ) . (3.4.1)

IRn IRn

() = i k () IRn , k = 1, . . . , n , (3.4.2)

xk

as it can be seen by applying an integration by parts in the integral which defines the Fourier

transform of /xk .

The following result gives the announced characterization of H1 (IRn ) in terms of summability

at infinity of the Fourier transform of its functions.

Proposition 3.4.1. One has

and Z 1/2

2 2

kvkH1 (IRn ) = (1 + kk )|v()| d .

IRn

v

Proof. It is enough to prove that, for each k = 1, . . . , n, x k

L2 (IRn ) in the distributional

sense iff the function 7 k v() belongs to L2 (IRn ), with identical norm. Let us assume that

v 2 n n

xk = gk L (IR ); then, using (3.4.2) and (3.4.1), for all D(IR ) one has on the one side

Z Z Z

v

h , i = v(x ) (x ) dx = v() ()d = i k v() () d

xk IRn xk IRn xk IRn

Z Z

v

h , i = gk (x )(x ) dx = gk ()() d .

xk IRn IRn

By equating the two last expressions and by recalling that is arbitrary, we get i k v() = gk ()

almost everywhere in IRn , and therefore the function 7 k v() belongs to L2 (IRn ). Conversely,

if this happens, one sets gk () = i k v(), so that its inverse transform gk (x ) belongs to L2 (IRn )

v

and satisfies gk = x k

.

The argument given in the proof shows that

1 n v

for all v H (IR ) , () = i k v() IRn , k = 1, . . . , n . (3.4.3)

xk

3.4. THE SPACES HS (IRN ) 45

Given a vector = (1 , . . . , n ) IRn and a multi-index = (1 , . . . , n ) INn , let us set

= 11 22 nn IR. An argument similar to the proof of Proposition 3.4.1 shows that

Dv L2 (IRn ) in the sense of distributions if and only if v() L2 (IRn ); in this case, one has

(D v) () = i|| v() ,

which generalizes (3.4.3). Thus, any Sobolev space Hm (IRn ) can be characterized as the subset

of L2 (IRn ) of those functions satisfying v() L2 (IRn ) for all INn such that || m; the

norm in Hm (IRn ) could be represented as

1/2

Z X

kvkHm (IRn ) = ||2|v()|2 d ,

IRn ||m

where || = (|1 |, |2 |, . . . , |n |). An equivalent but simpler expression of the norm is preferred,

which can be derived by applying the following technical lemma, whose elementary proof is left to

the reader.

Lemma 3.4.2. There exists constants c, C > 0 depending only on n and m such that

m X m

c 1 + kk2 ||2 C 1 + kk2 , IRn .

||m

The result allows us to characterize Hm (IRn ) in an equivalent manner as the subset of L2 (IRn )

m/2

of those functions such that 1 + kk2 v() L2 (IRn ), and to use the L2 (IRn )-norm of this

m n

function as an equivalent norm in H (IR ).

At this point, a remarkable observation can be made, namely, that the latter characterization

does not require the parameter m to be an integer: any real value of m is admissible. This leads

us to extend the definition of Sobolev spaces given so far, to the case of real positive indices.

s/2

Hs (IRn ) = {v L2 (IRn ) : 1 + kk2 v() L2 (IRn )}

Z 1/2

2 s 2

kvkHs (IRn ) = 1 + kk |v()| d . (3.4.4)

IRn

In this way, we obtain a continuos family of separable Hilbert spaces, which satisfy

s n s n

H (IR ) H (IR ) iff s >s ,

s n m n

H (IR ) = H (IR ) if s=m,

m+1

H (IR ) H (IR ) H (IRn )

n s n m

iff m<s<m+1.

The last relation shows that the Sobolev spaces Hs (IRn ) of non-integer index can be viewed as a

sort of interpolating spaces between consecutive Sobolev spaces of integer index. The concept can

be made rigorous, within the so-called Theory of Space Interpolation.

46 CHAPTER 3. SOBOLEV SPACES

One can furtherly extend the definition of Sobolev space to the case of negative indices, by

setting

Hs (IRn ) = H|s| (IRn ) if s < 0 ,

where X denotes the dual space of the Hilbert space X; equivalently, Hs (IRn ) can be defined

as the space of the distributions whose Fourier transform (defined in a suitable sense) makes the

right-hand side of (3.4.4) finite.

The idea is to apply to any function in H1 (IRn ) a truncation, which yields a compactly supported

function, followed by a regularization, which produces a C function.

Truncation. Given any R > 0, let R D(IR) be an even function satisfying

0 R (t) 1 t IR ,

R (t) 1 if |t| R ,

R (t) 0 if |t| R + 1 .

Then, given any v H1 (IRn ), one can prove that the function vR (x ) = R (kx k)v(x ) belongs to

H1 (IRn ) and is supported in B(0, R + 1); furthermore, kv vR kH1 (IRn ) 0 as R +.

Regularization. Given any > 0, let (x ) be any non-negative function in D(IRn ) satisfying

Z

supp B(0, ) , (x ) dx = 1 .

IRn

An example of such function is obtained by properly scaling the function given in Example 1.2.1.

Note that as 0, converges in D (IRn ) to the distribution 0 .

Then, given any v H1 (IRn ), one can prove that the convolution function

Z

v (x ) = ( v)(x ) = (x y )v(y ) dy

IRn

belongs to H1 (IRn ) and is infinitely differentiable at every x IRn ; furthermore, kvv kH1 (IRn ) 0

as 0.

Finally, we combine the two previous approximations by considering functions vR, = (vR )

obtained by first truncating a function v H1 (IRn ), and then regularizing the result. Since both

vR and are compactly supported, so is vR, ; precisely, supp vR, B(0, R + 1 + ). Thus, vR,

belongs to D(IRn ).

An appropriate choice of = (R), such that (R) 0 as R , shows that v can be

approximated in H1 (IRn ) to any prescribed precision by a function vR, for a sufficiently large R.

One of the simplest examples of open domain with nonempty boundary is the semi-space =

IRn+ = {x = (x1 , . . . , xn1 , xn ) IRn : xn > 0}, whose boundary IRn+ = {x = (x1 , . . . , xn1 , 0)

IRn } can be identified with IRn1 . For notational simplicity, every point x IRn will be written

as x = (x , xn ) with x IRn1 .

3.5. THE SPACE H1 (IRN

+) 47

A number of properties of H1 (IRn+ ), such as Property 3.2.8, can be obtained from the analogous

properties of H1 (IRn ) after introducing a suitable prolongation operator which extends the func-

tions belonging to H1 (IRn+ ) into functions belonging to H1 (IRn ). Precisely, given any v H1 (IRn+ ),

let us set

(P v)(x ) = v(x , |xn |) x = (x , xn ) IRn .

Thus, P realizes an extension of v by an even reflection around the boundary IRn+ . It is easy to

check that P v H1 (IRn ) iff v H1 (IRn+ ), and that

kP vkH1 (IRn ) 2kvkH1 (IRn+ v H1 (IRn+ ) ,

Property 3.5.1. There exists a linear continuous operator P : H1 (IRn+ ) H1 (IRn ) such that

P v|IRn+ v for all v H1 (IRn+ ).

Any operator satisfying the conditions of the Property is termed a prolongation operator in

H1 (IRn+ ).

As an application, we can easily prove Property 3.2.8 in the current situation. Indeed, given

any v H1 (IRn+ ), by the analogous result in IRn one can find a sequence of functions vk D(IRn )

which converge to v = P v in the H1 (IRn )-norm. It is immediate that the functions vk = (vk )|IRn

+

belong to D(IRn+ ) and converge to v in the H1 (IRn+ )-norm.

Next, we begin the discussion of the concept of trace on the boundary of a function

belonging to H1 (). Clarifying this concept is very important, for instance in order to give the

proper meaning to a Dirichlet boundary condition. The present geometrically simple situation

will allow us to keep ideas separated from technicalities, which may occur in the case of general

domains. While for a smooth function defined in (e.g., a function in D()) its trace on is

defined pointwise in the obvious way, for a function v which merely belongs to H1 () the same

procedure cannot be applied in dimension n 2. Indeed, v is an L2 -function, hence it is actually

a class of equivalence of functions, which may arbitrarily differ on subsets of zero measure in ;

since is precisely one of such subsets, the pointwise restriction of v to is meaningless. (The

situation is different in dimension 1, since we have seen that in each class of equivalence there is

one member which is continuous up to the boundary (Property 3.2.7), so that in particular its

boundary values are well-defined.)

The correct approach to the problem of defining boundary traces of functions in H1 () consists

of considering the trace operator as a linear continuous mapping defined on H1 (), which is first

defined pointwise on the subset of smooth functions, and which is next extended in a unique way

to the whole space thanks to the density of smooth functions.

We will detail this procedure in the case of = IRn+ . To this end, the following result is of

paramount importance.

Proposition 3.5.2. One has

k|IRn+ kL2 (IRn+ ) kkH1 (IRn+ ) D(IRn+ ) ,

Proof. Given any D(IRn+ ), let A > 0 be such that supp IRn1 [0, A]. For any x IRn1 ,

one has by the fundamental theorem of integral calculus in one dimension,

Z A Z A

2 2 2

2

(x , 0) = (x , A) (x , 0) = (x , xn ) dxn = 2 (x , xn ) dxn ;

0 xn 0 xn

48 CHAPTER 3. SOBOLEV SPACES

Z 1/2 !1/2

A Z A 2

2 2

(x , 0) 2 (x , xn ) dxn (x , xn ) dxn

0 0 xn

Z A Z A 2

2 (x , xn ) dxn + (x , xn ) dxn .

0 0 xn

Z Z Z 2

2 2

(x ) dx (x ) dx + (x ) dx , (3.5.1)

IRn

+ IRn

+ IRn

+

xn

The result tells us that the trace operator, which maps D(IRn+ ) into its restriction | on

, is continuous if the domain of definition is equipped with the H1 (IRn+ )-norm and the image

is equipped with the L2 (IRn+ )-norm. But since D(IRn+ ) is dense in H1 (IRn+ ), the Continuous

Extension Theorem guarantees that there exists a unique extension, say , which is defined on

the whole of H1 (IRn+ ), takes its values in L2 (IRn+ ) and is linear and continuous between these two

spaces. This means that the trace (v) of any function v H1 (IRn+ ) is well-defined as an element

of L2 (IRn+ ) (at least).

With a little additional effort, we can see that the image of is actually a proper subspace

of L2 (). Note, indeed, that in the proof of Proposition 3.5.2 only the L2 -norms of and x n

2

have been used on the right-hand side (see (3.5.1)). We now involve the L -norms of the other

first-order partial derivatives of with the aim of improving the bound given in the previous

proposition, by putting a stronger norm on the left-hand side. To this end, we recall that IRn+

is isomorphic to IRn1 , so that one can define the space H1/2 (IRn+ ) by identifying it to the space

H1/2 (IRn1 ) defined in Def. 3.4.3. The new result is as follows.

+

respect to the variables x , keeping xn fixed. Using the same notation as in the proof of Prop.

3.5.2, we have (x , xn ) 0 for all xn A, hence also ( , xn ) 0 for all xn A. Thus, we can

write

Z A Z A

2 2

| ( , 0)| = | ( , xn )| dxn = 2 IRe ( , xn ) ( , xn ) dxn

0 xn 0 xn

Z A

2

|( , xn )| ( , xn ) dxn .

0 xn

Z Z

2 1/2 2

(1 + k k ) | ( , 0)| d 2 (1 + k k ) |( , xn )|

2 1/2

( , xn ) d dxn

IRn1 IRn

+

xn

Z Z 2

2 2

(1 + k k )|( , xn )| d dxn + ( , xn ) d dxn .

xn

IRn+ IRn

+

3.6. SOBOLEV SPACES ON BOUNDED DOMAINS 49

Recalling Proposition 3.4.1, the first integral on the right-hand side equals

Z " n1 2 #

X

2 + (x ) dx ;

n

IR+ xk

k=1

on the other hand, the second integral on the right-hand side equals

Z

2

(x ) dx ,

IRn

+

xn

transform

in the x -variables are independent

of each other, so that they commute, i.e., x n

= x n

. This concludes the proof.

The previous result tells us that the trace operator maps H1 (IRn+ ) into H1/2 (IRn+ ) in a

continuous way. A deeper result, that we will not prove, guarantees that the image of H1 (IRn+ )

under is precisely the space H1/2 (IRn+ ), and that admits a continuous right-inverse :

H1/2 (IRn+ ) H1 (IRn+ ), i.e., ( (g)) = g for all g H1/2 (IRn+ ), or, equivalently,

(vg ) = g and kvg kH1 (IRn+ ) CkgkH1/2 (IRn ) .

+

termed trace operator, such that () = | for all D(IRn+ ). It is surjective upon

H1/2 (IRn+ ) and admits a continuous right-inverse

From now on, we suppose that is a bounded domain, and we make some assumptions on its

boundary , which will guarantee the validity of such results as the existence of prolongation or

trace operators, as for the half-space.

The following definitions will be crucial for our purposes.

Definition 3.6.1. A bounded open domain is said of class Cm (or simply a Cm -domain) for m 1

if there exists a finite covering of by open bounded sets Ai , i = 0, 1, . . . , I, such that

i) A0 A0 ;

ii) for each i = 1, . . . , I, there exists a mapping i : Ai B(0, 1) with the following properties:

a) i is bijective ;

b) i is of class Cm , with inverse i1 also of class Cm ;

50 CHAPTER 3. SOBOLEV SPACES

d) i (Ai ) = {y B(0, 1) : yn = 0} .

polyhedral (in dimension 3) type. The presence of corners or edges prevents them from being of

class C1 . The following definition relaxes the previous one for m = 1, in such a way that polygonal

or polyhedral domains fulfil its conditions.

Definition 3.6.2. A bounded open domain is said a Lipschitz domain if there exists a finite

covering of by open bounded sets Ai , i = 0, 1, . . . , I, such that

i) A0 A0 ;

ii) for each i = 1, . . . , I, there exists a mapping i : Ai B(0, 1) with the following properties:

a) i is bijective ;

b) i is of class C1 , with inverse i1 also of class C1 ;

c) there exists a Lipschitz-continuous function gi : IRn1 IR such that gi (0) = 0,

and

i (Ai ) = {y = (y , yn ) B(0, 1) : yn = g(y )} .

Obviously, a C1 -domain is a particular case of Lipschitz domain, where one can take each gi 0.

The concept of C1 -domain (of Lipschitz domain, resp.) can be equivalently expressed by saying

that locally its boundary is a graph of a C1 -function (a Lipschitz-continuous function, resp.), such

that the domain lies on one side of the graph.

One can prove that any convex domain is a Lipschitz domain.

Examples of domain which are neither C1 nor Lipschitz are those containing cusp points, such

as

= {x IR2 : x2 + (y 1)2 < 1 and x2 + (2y 1)2 > 1} ,

i.e., the region between two circumferences which are tangent at the origin. The origin is a cusp

point, such that in none of its neighborhoods the boundary can be represented as the graph

of a function.

The concept of partition of unity provides the tool which allows us to localize the study of a

function defined in .

Definition 3.6.3. A partition of unity associated with the covering {Ai }i=0,1,...,I of is a set of

nonnegative C -functions i : IRn IR such that

i) supp i Ai ;

I

X

ii) i (x) = 1 x .

i=0

Property 3.6.4. Given any finite covering of , there exists a partition of unity associated with

it.

3.6. SOBOLEV SPACES ON BOUNDED DOMAINS 51

Example 3.6.5. Let us exhibit a simple partition of unity associated with a covering of an interval

of the real line. Let us start by considering the even function

(

exp t211 |t| < 1 ,

(t) =

0 |t| 1 ;

(t) dt

(s) = R

+ ,

(t) dt

which satisfies (s) 0 for s 1, (s) 1 for s 1, (s) is strictly increasing in [1, 1]. Note

that 21 is an odd function, and this easily implies the identity (s) + (s) 1 in IR. Setting

(s) = (20s), we squeeze into the interval [1/20, 1/20] the transition region between the values

0 and 1.

Consider now the interval = (1, 1) and the covering

A0 = ( 43 + 10 1

), ( 43 + 10

1

) , A1 = 34 10 1 5 1

, 4 + 10 , A2 = ( 54 + 10 1

), 43 + 10

1

) .

0 (x) = (x+ 34 )+( 43 x)1 , 1 (x) = (x 43 )+( 54 x)1 , 2 (x) = (x+ 54 )+( 34 x)1 .

operators defined in H1 (), starting from the analogous operators in H1 (IRn+ ). Let us sketch the

idea. The starting point consists of writing a function v H1 () as

I

! I

X X

v(x ) = 1 v(x ) = i (x ) v(x ) = (i (x )v(x )) ,

i=0 i=0

P

i.e., setting vi (x ) = i (x )v(x ), we express v as v = Ii=0 vi , with supp vi Ai and kvi kH1 ()

CkvkH1 () for i = 0, . . . , I.

Now, v0 vanishes in a neighborhood of , hence, we can think of it as extended by zero outside

, i.e., v0 H1 (IRn ). On the other hand, for i = 1, . . . , I, we define vi = vi i1 , a function which

is supported in B(0, 1) IRn+ , so that it can be extended by zero to a function vi in IRn+ , satisfying

vi H1 (IRn+ ) with kvi kH1 (IRn+ ) Ckvi kH1 () .

Let P : H1 (IRn+ ) H1 (IRn ) be the prolongation operator defined in Sect. 3.5. Then, the

function

vi (x ) on Ai ,

vi (x ) = i (x ) (P vi ) i (x ) on Ai C ,

0 on IRn \ A , i

is an extension of vi which satisfies kvi kH1 (IRn ) Ckvi kH1 () . Finally, the global prolongation

P

operator is defined as P v = v0 + Ii=1 vi . Thus, we have established the following result.

Property 3.6.6. Let be a bounded Lipschitz domain. There exists a linear continuous operator

P : H1 () H1 (IRn ) such that P v| v for all v H1 ().

52 CHAPTER 3. SOBOLEV SPACES

As for the case = IRn+ , this property allows one to prove Property 3.2.8 for all Lipschitz domains.

Let us now consider the problem of defining the trace operator. To this end, let IRn+ :

H (IRn+ ) H1/2 (IRn+ ) be the trace operator defined in Sect. 3.5. Then, the function IRn+ (vi )i

1

zero to \ Ai , givingP rise to a function (vi ) defined on . Thus, the trace operator on

is defined as (v) = Ii=1 (vi ). It is easily seen that () = | if C (), and that

(v) L2 () with

k (v)kL2 () CkvkH1 () . (3.6.1)

its image, which is a subspace of L2 (). On the other hand, is clearly non-injective (many

functions in may have the same trace on ). Thus, let us introduce the subspace

Next, let us give a few results about quotient spaces. If X is a Hilbert space, with inner

product (x, y)X and norm kxkX , and X0 is a closed subspace of X, the quotient space X/X0 is

the set of all equivalence classes x = x + X0 = {y X : y x X0 }; it is a Hilbert space for the

quotient norm

kxkX/X0 = inf kxkX = inf kx + x0 kX .

xx x0 X0

We observe that the infimum above is actually a minimum. Indeed, given x X/X0 , there exists

a unique element x X such that kx kX = kxkX/X0 . This result can be proven by taking any

element y x and setting x = y y, where is the orthogonal projection operator from X

upon X0 ; it is easily seen that x is independent of the particular choice of y, and satisfies the

conditions stated above. Equivalently, the linear continuous mapping x X 7 x X/X0 admits

a continuous right-inverse x X/X0 7 x X.

We apply these results to the quotient space H1 ()/H10 (), after observing that, by definition

of kernel of a linear operator, induces an algebraic isomorphism between H1 ()/H10 () and

H1/2 (). Therefore, we can equip the space H1/2 () by the quotient norm

where v is the unique equivalence class of all functions v H1 () satisfying (v) = g; equiva-

lently, we have

|k g |kH1/2 () = inf kvkH1 () , (3.6.5)

vH1 (), (v)=g

or

|k g |kH1/2 () = kvg kH1 () ,

where v = vg is the element in v of smallest H1 ()-norm. One can prove (exercise) that vg is the

unique solution of the elliptic problem

(

vg + vg = 0 in ,

(3.6.6)

vg = g on .

3.7. THE SPACE H10 () AND THE POINCARE-FRIEDRICHS INEQUALITY 53

Thus, for the norm just introduced in H1/2 (), the mapping g 7 vg is a continuous right-inverse

of the mapping v 7 (v).

One can give an intrinsic definition of the space H1/2 (), as one of the fractional order

Sobolev spaces Hs (), 0 < s < 1, defined as

Z Z

s 2 2 |g(x ) g(y )|2

H () = {g L () : |g|Hs () = 2s+n

dxdy < +}

kx y k

1/2

kgkHs () = kgk2L2 () + |g|2Hs () .

This norm is equivalent to the norm defined in (3.6.4). If = IRn+ , this definition is also equivalent

to the one given in Sect. 3.5 via the Fourier transform.

We summarize the results that can be proven about the trace operator .

Theorem 3.6.7. Let be a bounded Lipschitz domain. There exists a linear continuous operator

: H1 () H1/2 () ,

termed trace operator, such that () = | for all C (). It is surjective upon H1/2 ()

and admits a continuous right-inverse

: H1/2 () H1 () ,

Remark 3.6.8. Let us show on a single example how the previous results can be extended to

Sobolev spaces of higher order.

Consider a bounded C2 -domain, and take a function v H2 (). Then, not only its trace (v)

is well-defined in H1/2 (), but also the traces (v/xi ) of its first-order partial derivatives

are well-defined in H1/2 ().

This is expressed by saying, on the one side, that (v) belongs to the Sobolev space H3/2 ()

1

(i.e., it is more regular, as a consequence of the fact that v is more regular

P than just an H ()-

function) and, on the other side, that the normal derivative v/n = i (v/xi )ni is well-

defined in H1/2 ().

The space H10 () has been defined in (3.6.3). It is the natural space in which to set the variational

formulation of a Dirichlet problem for a second-order elliptic boundary-value problem.

An equivalent definition is based on the following property.

Thus, H10 () can be equivalently defined as the closure of D() with respect to the topology

of H1 (); i.e., a function v H1 () belongs to H10 () if and only if there exists a sequence of

functions n D() satisfying kv n kH1 () 0 as n .

As a closed subspace of the Hilbert space H1 (), H10 () is itself a Hilbert space, for the

same inner product as in H1 (). On the other hand, a simpler, yet equivalent inner product

54 CHAPTER 3. SOBOLEV SPACES

can be defined in H10 () (equivalent meaning that it induces an equivalent norm). In order to

motivate its definition, let us start from the observation that the norm in H1 (), given in Definition

3.2.1, depends on both the L2 ()-norm of the function and the L2 ()-norm of its gradient. In

H1 () functions exist, which have one of the two norms much larger that the other one. For

instance, highly oscillatory but bounded functions (such as v,k (x, y) = sin kx cos ky in the square

= (0, 2)2 ) may be arbitrarily small while their gradients may be arbitrarily large. On the

contrary, constant functions (such as vk (x, y) = k) may be arbitrarily large while their gradients

are identically zero.

However, suppose bounded and consider a function constrained to vanish on : then, it

can be large somewhere in the domain only if its gradient, too, is large somewhere. This intuitive

concept can be made rigorous through an important inequality, known as the Poincare-Friedrichs

inequality, which we now state.

Proposition 3.7.2. Let be a bounded domain. Then, there exists a constant CP > 0 such that

Any constant for which this inequality holds is referred to as a Poincare constant in the domain.

There exists a minimal value CP () of this constant, depending only on , which can be referred

to as the Poincare constant of the domain .

Proof. We follow the strategy of first proving the inequality for all functions in D(). Next, since

D() is dense in H10 () and both sides of the inequality depend continuously on the H1 ()-norm,

we can pass to the limit and extend the inequality to all functions in H10 ().

Since is bounded (in particular, in the xn -direction), there exist constants a < b such that

IRn1 [a, b]. Given any D(), let us extend it by zero outside ; then, for any x IRn1

and any xn [a, b], the fundamental theorem of Calculus yields

Z xn

(x , xn ) = (x , s) ds .

a xn

By the Cauchy-Schwarz inequality, we get

Z Z xn 1/2 Z !1/2

xn

xn 2

|(x , xn )| =

1

(x , s)ds 2

1 ds

xn (x , s) ds

a xn a a

Z b !1/2

2

(xn a) 1/2

xn (x , s) ds .

a

Z b Z b Z b Z b

2 1 2

2

(x , xn ) dxn (xn a) dxn

(x , s) ds = (b a)2 (x , s) ds .

a a a xn 2 a xn

Z Z 2

1

2

(x ) dx (b a)2 (x ) dx .

2

IRn1 [a,b] IRn1 [a,b] xn

Z Z 2

1

2

kkL2 () = 2

(x ) dx (b a) 2 (x ) dx 1 (b a)2 kk2 2 ;

2 (L ())n

xn 2

3.7. THE SPACE H10 () AND THE POINCARE-FRIEDRICHS INEQUALITY 55

thus, inequality (3.7.1) is established with CP = 21/2 (b a) for any D(). The existence of

a minimal value of the Poincare constant will be proven in Chap. 6.

The assumptions make above in order to obtain the Poincare-Friedrichs inequality ( bounded

and functions vanishing on the whole of ) are just one of the possible sets of assumptions which

guarantee the inequality to hold. Here are possible extensions:

The proof clearly indicates that need not be bounded, but just bounded in one direction,

i.e., in the direction of a coordinate axis (possibly after a rigid rotation).

Functions need not vanish on the whole of , but just on a proper subset which has

positive (n 1)-dimensional measure, provided any point in the domain can be connected

to by a curve completely contained in the domain. In particular, if we introduce the

closed subspace of H1 () of the functions vanishing on a subset of positive (n 1)-

dimensional measure, i.e., if we set

main.

Let us introduce the closed subspace of L2 () of the zero-average functions, i.e., let us set

Z

L20 () 2

= {v L () : v = 0} . (3.7.3)

Note that zero-average functions cannot be strictly positive or strictly negative throughout

the domain; hence, if in addition they are continuous, they necessarily vanish somewhere

in the domain. So, it is not unexpected that one can prove (see Exercise 3.2) that the

Poincare-Friedrichs inequality holds in the closed subspace of H1 () given by H1 () L20 ().

We are now ready to introduce, as announced, a new inner product in H10 () or, more generally,

in any subspace of H1 () for which a Poincare-Friedrichs inequality holds.

Proposition 3.7.3. Let H0 denote any closed subspace of H1 () for which there exists a constant

CP > 0 such that

kvkL2 () CP kvk(L2 ())n v H0 . (3.7.4)

(u, v)H0 = u v (3.7.5)

R

2 1/2

is an inner product in H0 ; the induced norm kvkH0 = kvk is equivalent to the standard

H1 ()-norm in H0 , since one has

1

2

kvkH1 () kvkH0 kvkH1 () v H0 . (3.7.6)

1+CP

Proof. It is enough to prove the first inequality in (3.7.6), which immediately follows from

(3.7.4).

56 CHAPTER 3. SOBOLEV SPACES

In this section, we present, without proof, two fundamental theorems concerning Sobolev spaces,

the Rellich theorem and the Sobolev imbedding theorem.

Rellichs theorem is the counterpart in Sobolev spaces of classical theorems such as the Ascoli-

Arzela theorem. It pertains to the possibility of extracting a converging sequence of functions,

with respect to a certain Sobolev norm, from knowing that their derivatives of sufficiently high

order are bounded in the L2 ()-norm. The essential condition is that the domain has to be

bounded, for otherwise the result is not true.

The precise statement is as follows.

Theorem 3.8.1. (Rellich) Let be a bounded domain in IRn . Then, for any m 0 the inclusion

Hm () Hk () with 0 k < m, is compact.

For instance, if {vn }n1 is a sequence of functions in H1 () satisfying kvn kH1 () C for

some constant C > 0, then the theorem assures the existence of a subsequence {vnj }j1 which is

convergent in L2 ().

The Sobolev imbedding theorem links Sobolev regularity to classical regularity, allowing one

e.g. to see the weak solution of a variational problem as a classical solution as well. In essence,

the theorem says that any function in Hm () for large enough m has the property that all its

derivatives of order up to a certain k < m are classical derivatives (and not just distributional

derivatives), and they are continuous in . Furthermore, even if m is not large enough, any

function in Hm () is p-integrable for some p > 2 (and not just square-integrable), as soon as

m > 0.

The precise statement is as follows.

Theorem 3.8.2. (Sobolev) Let be any domain in IRn . Let m > 0 be given, and denote by [z]

the largest integer z.

n2m > 2.

iii) If m > n/2, then Hm () Ck, () L (), for k = [m n/2] and = m n/2 k if

m n/2 is not an integer, and for k = m n/2 1 and arbitrary < 1 if m n/2 is an

integer.

All inclusions above are continuous. In addition, if is bounded, the inclusions are compact.

Examples 3.8.3.

i) In dimension n = 1, the theorem gives Hm () Cm1,1/2 (). This result, for m = 1, has

been already proven in Sect. 3.2 (Property 3.2.7).

ii) In dimension n = 2, the theorem gives H1 () Lp () for all p < (but not contained in

L (), as already noted in Sect. 3.2), and Hm () Cm2, () L () for any integer m 2

iii) In dimension n = 3, the theorem gives H1 () L6 (), and Hm () Cm2,1/2 ()L ()

for any integer m 2.

3.9. THE DUALS OF H1 () AND H10 () 57

In this section, we analyze the structure of the dual spaces of the Sobolev spaces H1 () and H10 ().

The results will tell us which kind of right-hand side is admissible in the variational formulation

of a second-order elliptic problem.

We recall that the dual space of a Banach space X is the Banach space X of the continuous

linear forms F : X IR, equipped with the norm

F (x)

kF kX = sup .

xX kxkX

If X is a Hilbert space, the Riesz Representation Theorem says that each F X can be written

as F (x) = (y, x)X x X, for a unique y X, which satisfies kykX = kF kX . Thus, X can be

identified to X via the isometry F 7 y.

If X and Y are Banach spaces satisfying X Y with continuous injection, i.e., kxkY CkxkX

for all x X, then any FY Y defines an FX X by setting FX (x) = FY (x) x X;

furthermore, kFX kX CkFY kY , i.e., the mapping FY Y 7 FX X is continuous. If,

in addition, X is dense in Y , then this mapping is injective, since FX = GX means FY (x) =

GY (x) x X, whence FY (y) = GY (y) y Y , by the uniqueness of the extension of a continuous

map defined on a dense subset. In other words, we can identify Y to a subspace of X , i.e., we

have Y X with continuous injection and dense image.

A particularly important situation is the following one. We are given two Hilbert spaces V

and H, such that V H with continuous injection and dense image; furthermore, we identify H

to H via the Riesz Representation Theorem as above. Then, we have H = H V , so that we

can write the chain of inclusions

V H V , (3.9.1)

where each inclusion is continuous and with dense image. The pair (V, H) is often called a Gelfand

pair; equivalently, the triple (V, H, V ) is called a Gelfand triple. Examples of Gelfand triples are

(H1 (), L2 (), (H1 ()) ) and (H10 (), L2 (), (H10 ()) ) ,

arbitrarily well by compactly supported, smooth functions, i.e., D() is dense in L2 (), so that

a fortiori both H10 () and H1 () are dense in L2 ().

We want to find an explicit representation of any element in (H1 ()) or (H10 ()) . To this end,

let us introduce the mapping

S : H1 () (L2 ())n+1

v v (3.9.2)

v 7 v, , . . . , = Sv .

x1 xn

The mapping is trivially injective, since Sv = 0 implies that the first component of Sv, i.e., v

itself, is zero; furthermore, it is an isometry, since kSvk(L2 ())n+1 = kvkH1 () for all v H1 ().

Thus, we can identify H1 () to the subspace Z = S(H1 ()) of (L2 ())n+1 ; this subspace is closed,

since H1 () is complete. Correspondingly, the dual of H1 () can be identified to the dual of Z,

via the isometry F (H1 ()) 7 FZ Z defined as FZ (w) = F (S 1 (w)) w Z.

We now invoke the Hahn-Banch Theorem, which guarantees that, given a linear continuous

form FZ on a closed subspace Z of a Banach space X, there exists a linear continuous form Fe on

58 CHAPTER 3. SOBOLEV SPACES

X which extends FZ and such that kFZ kZ = kFekX . Thus, if we start from any F (H1 ()) ,

there exists Fe (L2 ())n+1 ) such that

and kF k(H1 ()) = kFZ kZ = kFek(L2 ())n+1 ) . On the other hand, having identified (L2 ()) to

L2 (), the space (L2 ())n+1 ) can be identified to (L2 ())n+1 , so that Fe can be identified to an

element f = (f0 , f1 , . . . , fn ) (L2 ())n+1 by the relation

n

X

Fe(w) = (f , w)L2 ()n+1 = (fi , wi )L2 () w (L2 ())n+1 , (3.9.4)

i=0

P 1/2

and kFek((L2 ())n+1 ) = kf k(L2 ())n+1 = n 2

i=0 kfi kL2 () . Combining (3.9.3) and (3.9.4), we

conclude that F (H1 ()) can be represented as

n

X

v

F (v) = (f0 , v)L2 () + fi , v H1 () , (3.9.5)

xi L2 ()

i=1

P 1/2

n

for suitable functions f0 , f1 , . . . , fn L2 () satisfying kF k(H1 ()) = 2

i=0 kfi kL2 () .

Since the extension FZ 7 Fe is not unique, the representation (3.9.5) is not unique as well.

For instance, since Z Z

v

dx = v dx if D() ,

x1 x1

we can replace in (3.9.5) f1 by f1 + and f0 by f0 x 1

without changing the right-hand side. If

2 n+1

f = (f0 , f1 , . . . , fn ) (L ()) denotes now any (n + 1)-ple of functions for which (3.9.5) holds,

then by the Cauchy-Schwarz inequality we have

n

X
v

|F (v)| kf0 kL2 () kvkL2 () +

kf1 kL2 ()
kf k(L2 ())n+1 kvkH1 () ,

xi
L2 ()

i=1

i.e., kF k(H1 ()) kf k(L2 ())n+1 . On the other hand, the construction above shows that there

exists f (L2 ())n+1 for which the equality sign is attained.

The following statement summarizes the results obtained so far.

n

X

v

F (v) = (f0 , v)L2 () + fi , v H1 () , (3.9.6)

xi L2 ()

i=1

(3.9.6) holds }, we have

n

!1/2

X

2

kF k(H1 ()) = min kf k(L2 ())n+1 = min kfi kL2 () . (3.9.7)

fR(F ) fR(F )

i=0

3.9. THE DUALS OF H1 () AND H10 () 59

We now consider the dual space of H10 (), which is usually denoted by H1 (). All previous

considerations can be repeated, with the only change that the operator S introduced in (3.9.2)

is now restricted to H10 (), and consequently its image is a closed subspace of Z, say Z0 . Then,

given any F H1 (), we arrive as above to the representation formula

n

X

v

F (v) = (f0 , v)L2 () + fi , v H10 () , (3.9.8)

xi L2 ()

i=1

The key difference with respect to case discussed above is that elements in H1 () are distri-

butions. Indeed, D() is dense in H10 (), hence, H1 () can be identified to a dense subspace

of D (), the space of distributions over . In other words, we have the sequence of continuous

inclusions with dense images

D() H10 () L2 () H1 () D () .

Recalling the expression of the partial derivatives of a distribution (Def. 2.2.1), the right-hand

side of (3.9.8) can be written as

n * n

+

X X fi

(f0 , )L2 () + fi , = f0 , D() ,

xi L2 () xi

i=1 i=1

i.e.,

n

X fi

F = f0 = f0 (f1 , . . . , fn ) in D () . (3.9.9)

xi

i=1

Thus, any F H1 () is the sum of an L2 ()-function and the divergence (in the sense of

distributions) of a vector of L2 ()-functions. The representation of its norm is analogous to the

one in (H1 ()) .

We summarize the results as follows.

Theorem 3.9.2. Let us denote by H1 () the dual space of H10 (). Then, H1 () D (), and

any F H1 () can be represented, in a non-unique way, as in (3.9.9) for suitable f0 , f1 , . . . , fn in

L2 (); equivalently, (3.9.8) holds. In addition, setting R0 (F ) = {f (L2 ())n+1 : (3.9.8) holds },

we have !1/2

Xn

kF kH1 () = min kf k(L2 ())n+1 = min kfi k2L2 () . (3.9.10)

fR0 (F ) fR0 (F )

i=0

Examples 3.9.3.

i) Let = (a, b) be an interval in IR and let x0 any point in . Since H10 () H1 () C0 (),

the linear form v 7 Fx0 (v) = v(x0 ) belongs to both H1 () and (H1 ()) . As an element of

H1 (), it coincides with the distribution x0 (we simply say that x0 belongs to H1 ()); indeed,

we can represent it as in (3.9.9) setting Fx0 = f0 df

dx , with f0 0 and f1 (x) = H(x x0 ), where

1

H is the Heaviside function. We refer to Exercise 3.3 for a representation of Fx0 , as an element of

(H1 ()) , in the form (3.9.6).

ii) In dimension n > 1, the distribution x 0 , with x 0 , belongs neither to H1 () nor to

(H1 ()) ; indeed, neither H1 () nor H10 () are imbedded in C0 ().

60 CHAPTER 3. SOBOLEV SPACES

domain which is cut into two non-empty subsets and + according to where x1 < c or

x1 > c. Then, the form Z

F (v) = (v) d ,

where (v) is the trace of v on , belongs to both H1 () and (H1 ()) , since the mapping

: H1 ( ) L2 () is continuous, as seen in 1

R Sect. 3.6. As an element of H (), F coincides

with the distribution such that h , i = | d for all D(). It can be represented in

the form (3.9.9) by setting f0 = f2 = = fn 0 and f1 = + , the characteristic function of

the set + . We refer to Exercise 3.4 for a representation of F , as an element of (H1 ()) , in the

form (3.9.6).

iii) Let Lw = (Aw) + a w + a0 w be any second-order operator in ; let us assume

that all its coefficients belong to L (). Then, for any w H1 (), one has Aw (L2 ())n as

well as a w + a0 w L2 (). Thus, according to Thm. 3.9.2, Lw belongs to H1 () and the

mapping w H1 () 7 Lw H1 () is continuous. A particularly relevant case occurs when

Lw = w: by restricting w to H10 (), we obtain that the operator maps H10 () into its dual

H1 (). Precisely, there exists a constant C > 0 depending on such that

The Laplacian is actually an isomorphism between these two spaces, as discussed in the next

chapter (see Property 4.3.7).

One can also consider Lw, for w H1 (), as an element of (H1 ()) , by setting

Z Z

F (v) = (Aw) v + (a w + a0 w) v v H1 () .

iv) Let be a bounded Lipschitz domain; recall that the trace operator is continuous

from H1 () to L2 (). Thus, given any h L2 (), the form Fh defined by

Z

v 7 Fh (v) = h (v) d (3.9.12)

belongs to (H1 ()) . How can we represent Fh according to (3.9.6)? To answer this question, let

us consider the Neumann problem

w + w = 0 in ,

w = h on .

n

Anticipating the results of next chapter, we can say that there exists a unique solution w H1 ()

of the variational formulation of this problem, which is given by

Z Z

w v + wv = h (v) v H1 () . (3.9.13)

w w

Then, the (n + 1)-ple f = w, x 1

, . . . , xn provides a representation of the form Fh . It is not

difficult to prove, indeed, that this is the representation which realizes the minimum in (3.9.7).

3.10. EXERCISES 61

Consider now the same mapping (3.9.12), but restrict it to the functions of H10 (). Obviously,

(Fh )|H10 () belongs to H1 (), but... it is nothing else than the null form, i.e., Fh0 = 0

Fh0 =

1

H (). Since, by (3.9.12),

Z Z

0 = Fh0 () = h | = w + w = hw + w, i D() ,

the same f as above does provide one of the possible representations of Fh0 in the form (3.9.9),

but surely it is not the one with minimal norm! Obviously, such a representation is provided by

the (n + 1)-ple f = (0, 0, . . . , 0).

3.10 Exercises

3.1. Let us consider the function

u(x, y) = |x y|

on the set = [1, 2] [1, 2] IR2 , where is a real parameter.

(i) Find the values of for which u L2 () and those for which u H1 ().

(ii) For the values of that allow u to be in H1 (), calculate the Laplacian u and find the

space which it belongs to.

Z

kvk = v(x ) dx + kvk(L2 ())n

is a norm in H1 (), which is equivalent to the standard norm kvkH1 () . Deduce from this the

validity of a Poincare-Friedrichs inequality in the space H1 () L20 (), where L20 () is defined in

(3.7.3).

3.3. Prove that the form Fx0 defined in Example 3.9.3, i) can be represented as

Z b Z b

Fx0 (v) = f0 (x)v(x) dx + f1 (x)v (x) dx v H1 (a, b) ,

a a

where ( (

0 in (a, x0 ) , 0 in (a, x0 ) ,

f0 (x) = f1 (x) =

w(x) in (x0 , b) , w (x) in (x0 , b) ,

and w is the solution of the Neumann problem in (x0 , b):

(

w + w = 0 in (x0 , b) ,

w (x0 ) = 1, w (b) = 0 .

3.4. Adapt the arguments of the previous exercise in order to find a representation of the form

F defined in Example 3.9.3, ii), as an element of (H1 ()) .

62 CHAPTER 3. SOBOLEV SPACES

Chapter 4

Elliptic Problems

Hereafter, we will study a more general situation than the homogeneous Dirichlet boundary-value

problem given in (3.1.1). Precisely, let us consider the mixed Dirichlet/Neumann problem

Lu = (Au) + (au) + a0 u = f in ,

u = g on D , (4.1.1)

u

= h on N .

nA

Here, is a bounded and connected Lipschitz domain, whose boundary is partitioned into two

relatively open subsets D (the Dirichlet part of the boundary) and N (the Neumann part), i.e.,

they satisfy

D , N , D N = , D N = .

We assume again that the coefficients A, a, a0 satisfy (3.1.12); in addition, if N is not empty, we

assume that A and a are defined and continuous in a neighborhood of N , so that the conormal

vector nA = An and the normal component an = a n of the vector a make sense therein. We

u

recall that the conormal derivative of u is defined as n A

= nA u; see also (3.1.5).

We now add the crucial assumption that the equation is elliptic throughout the domain, i.e.,

the matrix A(x) is positive-definite at each x .

As far as the data f , g and h are concerned, we assume as in (3.1.12) that f L2 (), that g

is the trace on of a function of H1 (), i.e., g H1/2 (), and finally that h L2 (N ).

We now formulate the problem in a weak sense, so that the solution u is only required to

belong to H1 (). We recall the identity (3.1.8), obtained after multiplying equation Lu = f by

a test function v, integrating over and performing two integrations-by-parts. If we restrict

test functions to those vanishing on N , and replace therein the conormal derivative of u by the

prescribed value h, we obtain

Z Z Z Z Z Z

(Au) v u a v + a0 uv + an uv = fv + hv . (4.1.2)

N N

Note that all integrals which appear in this equation make sense if u, v H1 (), thanks to the

definition of this space and the trace properties established in Chapter 3. Hence, we introduce the

bilinear form

Z Z Z Z

a(u, v) := (Au) v u a v + a0 uv + an uv (4.1.3)

N

63

64 CHAPTER 4. ELLIPTIC PROBLEMS

Z Z

F (v) := fv + hv (4.1.4)

N

defined in H1 (). We also introduce the closed subspace of H1 () of the functions vanishing on

D , i.e., we set

H10,D () = {v H1 () : v = 0 on D } ,

Then, we are led to considering the following weak formulation of the boundary-value problem

stated in (4.1.1):

a) Homogeneous Dirichlet condition, g = 0

Both the solution u and the test functions v belong to the same space H10,D (). The weak

formulation becomes:

This case can be reduced to the previous one by a change of unknown. Precisely, since g

H1/2 (), there exists ug H1 () such that ug | = g; we write u = u0 + ug , where the new

unknown w satisfies (u0 )|D = u|D (ug )|D = (g g)|D = 0, i.e., u0 H10,D (). Substituting

the expression of u into (4.1.5), we obtain

still defined in H1 (), we reduce the non-homogeneous problem for u to the following one for w:

This change of dependent variable is made in preparation of applying the abstract existence and

uniqueness theory developed in the subsequent section.

Obviously, if D = , we simply consider Problem 4.1.2 with H10,D () = H1 ().

4.1. WEAK FORMULATION OF ELLIPTIC BOUNDARY-VALUE PROBLEMS 65

Examples 4.1.4. Let us consider the homogeneous Dirichlet problem for the Poisson equation:

(

u = f in ,

(4.1.8)

u = 0 on .

It admits the following weak formulation: Find u H10 () such that

Z Z

u v = fv v H10 () . (4.1.9)

We observe that, keeping in mind Theorem 3.9.2, the data f L2 () can be replaced by a data

F H1 (), i.e., of the form

n

X fi

F = f0 , fi L2 () .

xi

i=1

In this case, (4.1.9) becomes

Z Z n Z

X v

u v = f0 v + fi v H10 () . (4.1.10)

xi

i=1

ii) Let us consider the non-homogeneous Neumann problem for the Helmholtz equation:

(

u + a0 u = f in ,

u

(4.1.11)

n = h on .

Z Z Z

(u v + a0 uv) = fv + hv v H1 () . (4.1.12)

In this case, too, a more general data F can replace f on the right-hand side.

iii) Let us consider the mixed Dirichlet/Neumann problem for the convection-diffusion equa-

tion:

u + (au) =

f in ,

u = g on D , (4.1.13)

u

n = 0 on N ,

where > 0 is a constant and D = {x : (a n)(x) < 0}. It admits the following weak

formulation: Find u H1 (), with u = g on D , such that

Z Z Z Z

u v u a v + an uv = fv v H10,D () . (4.1.14)

N

The reduction to an equivalent homogeneous problem is based on the representation u = u0 + ug ,

where ug is any fixed function in H1 () satisfying u = g on . Then, u0 is defined by the

following weak formulation: Find u0 H10,D () such that

Z Z Z

u0 v u0 a v + an u0 v (4.1.15)

N

Z Z Z Z

= fv ug v + ug a v an gv v H10,D () . (4.1.16)

N

Next section will be devoted to study an abstract form of Problem 4.1.2 or 4.1.3, providing a

set of assumptions which guarantee its solvability.

66 CHAPTER 4. ELLIPTIC PROBLEMS

Let V be a reflexive Banach space, with norm kvkV . (We recall that a Banach space is reflexive

if its bidual V = (V ) can be identified with V itself; each Hilbert space is reflexive, thanks to

Riesz Representation Theorem.)

Let

a: V V R

(4.2.1)

(w, v) 7 a(w, v)

be a bilinear form defined in V ; let

F : V R

(4.2.2)

v 7 F (v)

In order to prove the solvability of this problem, we assume that the forms a and F are

continuous. The latter condition is equivalent to F V ; concerning the form a, we give the

following definition.

Definition 4.2.2. The bilinear form a is said to be continuous in V if there exists a constant

C > 0 such that

|a(w, v)| C kwkV kvkV w, v V . (4.2.4)

a(w, v)

kak = sup , (4.2.5)

wV,vV kwkV kvkV

If the form a is continuous, we can associate to it, in a canonical way, a linear operator

A:V V ,

defined as follows: for any w V , the mapping v V 7 a(w, v) R is linear and continuous;

therefore, it is an element of V , which we denote by Aw. In other words, A is defined by the

relations

(Aw)(v) = a(w, v) w, v V . (4.2.6)

Au = F .

4.2. THE LAX-MILGRAM THEOREM 67

It will be convenient to introduce the duality pairing between V and V , i.e., the bilinear form

hAw, vi = a(w, v) w, v V ,

hAu, vi = hF, vi v V .

Next, we introduce the crucial property of the bilinear form a, which will provide, together

with continuity, a sufficient condition for the solvability of Problem 4.2.1.

Definition 4.2.4. The bilinear form a is said to be coercive in V if there exists a constant > 0,

such that

a(v, v) kvk2V v V . (4.2.8)

Any > 0 satisfying (4.2.8) is termed a coercivity constant of the form a. The best coercivity

constant is given by

a(v, v)

= inf . (4.2.9)

vV kvk2 V

We are now ready to state the main result of the present theory.

Theorem 4.2.5. (Lax-Milgram) Assume that the bilinear form a on the reflexive Banach space

V is continuous and coercive, with coercivity constant . Then, given any linear continuous form

F on V , Problem 4.2.1 admits one and only one solution, which satisfies

1

kukV kF kV . (4.2.10)

Remark 4.2.6. The Lax-Milgram Theorem provides sufficient conditions under which Problem

4.2.1 is well-posed in the sense of Hadamard. This means that:

for any data F , the problem has a solution u;

the solution is unique;

the solution depends on the data in a continuous way.

The latter condition is made explicit as follows: let F1 , F2 be two elements in V , and let u1 , u2 V

be the solutions of the corresponding problems. Then, by linearity, the difference u1 u2 is the

solution of

a(u1 u2 , v) = (F1 F2 )(v) v V ;

hence, by (4.2.10) applied to this problem, we conclude that

1

ku1 u2 kV kF1 F2 kV .

In other words, a small perturbation in the data yields a small perturbation in the solution.

Equivalently, the mapping F V 7 u = A1 F V is Lipschitz-continuous, with Lipschitz

constant 1/.

68 CHAPTER 4. ELLIPTIC PROBLEMS

Remark 4.2.7. The Lax-Milgram Theorem can be viewed as a generalization of the Riesz Rep-

resentation Theorem in a Hilbert space. Indeed, assume that V is such a space and that a(w, v)

is a continuous and coercive bilinear form on V , which in addition is symmetric, i.e.,

p the bilinear form (w, v)a = a(w, v) is an inner product in V , which

induces a norm, kvka = a(v, v), equivalent to the original norm in V ; indeed, by coercivity and

continuity one has

kvk2V a(v, v) kak kvk2V v V ,

i.e., p

kvkV kvka kak kvkV v V .

Thus, any linear form on V which is continuous in the original norm of V is also continuous in the

norm induced by a, and vice-versa. Given any such form F , the Lax-Milgram Theorem assures

the existence of a unique u V such that

Proof of Theorem 4.2.5 At first, let us remark that if we assume for the moment the existence

of a solution, then its uniqueness and the bound (4.2.10) are immediate. Indeed, this inequality

follows from coercivity after choosing v = u in (4.2.3, since

dividing by kukV , we get the result. Uniqueness follows from (4.2.10): if u1 , u2 are any two

solutions of Problem 4.2.1, their difference satisfies

Thus, we are left with the task of proving the existence of a solution of Problem 4.2.1. We

will actually prove that the operator A introduced in (4.2.6) is an (algebraic and topological)

isomorphism between V and V . Let us proceed in several steps.

Step 1): A : V V is continuous. This follows, as expected, from the continuity of a; indeed,

for any w V ,

hAw, vi a(w, v) kak kwkV kvkV

kAwkV = sup = sup sup = kak kwkV .

vV kvkV vV kvkV vV kvkV

Step 2): A is injective. This follows from the coercivity of a. Indeed, Aw = 0 V means

hAw, vi = a(w, v) = 0 for all v V ; taking v = w we get kwk2V a(w, w) = 0, which imples

w = 0.

Let us introduce the image of A in V , i.e., the subspace

Then, A is an algebraic isomorphism between V and Z, whose inverse will be denoted, as usual,

by A1 .

4.2. THE LAX-MILGRAM THEOREM 69

Step 3): A1 : Z V is continuous. This follows again from the coercivity of a. Indeed, given

any G Z, let w V be such that Aw = G. This means that w satisfies

a(w, v) = G(v) v V ;

1

applying the bound (4.2.10) to this problem, we get kA1 GkV kGkV ,

which is precisely the

claim.

Step 4): Z is a closed subspace of V . This follows from the completeness of V . Indeed, let F

belong to the closure of Z in V . Then, there exist a sequence {Gn }nN converging to F in the

norm of V . Let us set wn = A1 Gn V . By the result of Step 3), we have

1

kwn wm kV kGn Gm kV n, m N .

This implies that {wn }nN is a Cauchy sequence in V , hence, it is converging to some w V

since this space is complete. Then, the form G = Aw belongs to Z and, by Step 1), the sequence

{Gn = Awn }nN converges to G in V . Since it also converges to F by assumption, necessarily

F = G by the uniqueness of the limit, hence, F belongs to Z.

Step 5): Z coincides with V . This follows from the reflexivity of V . By contradiction, assume

that Z is a proper subspace of V . Then, according to the Hahn-Banach theorem, there exists

a linear continuous form W on V such that W(G) = 0 for all G Z, but W(F ) 6= 0 for some

F V \ Z. In other words, W V , and kWkV > 0.

Since V is reflexive, W can be identified with an element w in V ; precisely, there exists a

unique w V such that

kwk2V a(w, w) = 0 ,

which implies kwkV = 0, i.e., w = 0. This contradicts the property kwkV > 0 stated before, and

the claim is proven.

The proof of the Lax-Milgram Theorem is then concluded.

Let us consider again the relevant case of a continuous and coercive bilinear form a(w, v), which

is symmetric, i.e. it satisfies (4.2.11); let us show that in this case Problem 4.2.1 is equivalent to

a minimization problem in V . Precisely, let us introduce the functional

J: V R

(4.2.12)

v 7 J(v) = 21 a(v, v) F (v) ,

70 CHAPTER 4. ELLIPTIC PROBLEMS

vV

Property 4.2.9. For any w V , the following identity holds:

J(w + v) = J(w) + a(w, v) F (v) + 21 a(v, v) v V . (4.2.14)

Proof. The result is an immediate consequence of the bilinearity of a and the linearity of F .

The identity can be interpreted as follows: think of v as an increment v = w given to w. Then,

w 7 a(w, w)F (w) represents the linear part of the increment of J, whereas w 7 12 a(w, w)

represents the quadratic part. Equivalently stated, (4.2.14) is nothing but the Taylor expansion

of J around w

h(HJ) v1 , v2 i = a(v1 , v2 ) v1 , v2 V .

Property 4.2.10. The functional J is strictly convex in V , i.e., it satisfies

J(v1 + (v2 v1 )) < J(v1 ) + J(v2 ) J(v1 ) .

J(v1 + (v2 v1 )) = J(v1 ) + a(v1 , v2 v1 ) F (v2 v1 ) + 12 2 a(v2 v1 , v2 v1 ) . (4.2.17)

On the other hand, using the symmetry of the bilinear form, we obtain the relation

J(v1 + (v2 v1 )) = J(v1 ) + J(v2 ) J(v1 ) 12 (1 )a(v2 v1 , v2 v1 ) .

gives the result.

The last property concerns the behavior of F at infinity.

4.2. THE LAX-MILGRAM THEOREM 71

1 2 2 1 kF kV

J(v) 2 kvkV kF kV kvkV = kvkV 2 .

kvkV

4, whence

J(v) 14 kvk2V ,

Properties 4.2.10 and 4.2.11 should give the intuitive idea that the graph of J in V R behaves

like an elliptic paraboloid in R3 (but beware: here we are in infinite dimension!). In particular,

the existence of a unique minimum for J should not be a surprise at this point. This is made

precise in the following fundamental statement.

Theorem 4.2.12. The weak Problem 4.2.1 and the minimization Problem 4.2.8 are equivalent:

u is a solution of the latter problem if and only if it is a solution of the former problem.

Thus, since we already know that Problem 4.2.1 admits one and only one solution, the same

is true for Problem 4.2.8.

Proof. At first, assume that u is a solution of Problem 4.2.1. Property 4.2.9 with w = u yields

6 0, then a(v, v) > 0 by coercivity, hence J(u + v) > J(u), i.e., u is a strict minimizer of J.

If v =

Conversely, let u be a solution of Problem 4.2.8. For any fixed v V , consider the quadratic

function (parabola) : R R defined by

() = J(u + v) = J(u) + a(u, v) F (v) + 21 2 a(v, v)

d

0= (0) = a(u, v) F (v) ;

d

since v is arbitrary, this shows that u is a solution of Problem 4.2.1.

In view of (4.2.15), equations (4.2.3) state that J(u) = 0, i.e., they express the property that

the functional J is stationary at u. They are often called the Euler-Lagrange equations of the

minimization problem.

If we consider the bilinear form a(w, v) defined in (4.1.3), then it is symmetric if (and only if)

A is a symmetric matrix and a = 0 throughout the domain. In this case, the weak formulation

of the boundary-value problem (4.1.1), given by (4.1.5), may be referred to as the variational

formulation of the problem, and the corresponding solution u is called the variational solution.

Example 4.2.13. Consider the homogeneous Dirichlet problem for the Poisson equation, given

in (4.1.8), and the related variational formulation (4.1.9). The corresponding functional J :

H10 () R is given by Z Z

1 2

J(v) = kvkRn fv .

2

72 CHAPTER 4. ELLIPTIC PROBLEMS

The first integral on the right-hand side is called the Dirichlet integral of v in .

In a (extremely simplified) description of small deformations in linear Elasticity (such as in

the membrane problem), v represents an admissible displacement, constrained to vanish on the

boundary of the body which occupies the domain . Then, the quantity

Z

1

kvk2Rn

2

represents the internal elastic energy associated with the displaced configuration, whereas the

integral

Z

fv

represents the potential energy associated with the work of the external force of density f when

the displacement v takes place. Thus, J(v) represents the total energy of the configuration de-

scribed by v. Eq. (4.2.13) translates the well-known physical principle that, among all admissible

displacements, the one which corresponds to the equilibrium of the elastic body under the external

forces is characterized by having the minimal total energy.

of elliptic boundary-value problems

Let us go back to the mixed boundary-value problem (4.1.1), and let us establish suitable assump-

tions under which Problems 4.1.2 or 4.1.3 are well-posed. Thanks to the Lax-Milgram theorem,

this will be accomplished by checking that the bilinear form a(w, v) defined in (4.1.3) is continuous

and coercive in H10,D (), and that the linear forms F (v) or F (v), defined in (4.1.7) or (4.1.4), are

continuous in H10,D ().

At first, let us deal with continuity.

Lemma 4.3.1. Under the assumptions on the coefficients A, a, a0 stated at the beginning of

Sect. 4.1, one has

where C(A, a, a0 ) depends upon kAk(L ())nn , kak(L ())n , ka0 kL () and ka nkL (N ) .

Proof. In the proof, C will denote a constant, independent of w, v and the coefficients of the

operator, which be different from place to place.

Using Holders inequality (3.1.11) with p = , p = p = 2, we easily bound each addend in

the definition of a(w, v); precisely:

Z

(Aw) v CkAk(L ())nn kwk 2

(L ())n kvk(L2 ())n

CkAk(L ())nn kwkH1 () kvkH1 () ,

Z

w a v kak(L ())n kwk 2 kvk 2

L () (L ())n kak(L ())n kwkH1 () kvkH 1 () ,

4.3. WELL POSEDNESS OF ELLIPTIC PROBLEMS 73

Z

a0 wv ka0 kL () kwk 2 kvk 2

L () L () ka0 kL () kwkH 1 () kvkH1 () ,

Z

an wv ka nkL (N ) kwkL2 (N ) kvkL2 (N ) ka nkL (N ) kwkL2 () kvkL2 ()

N

Cka nkL (N ) kwkH1 () kvkL2 (H1 ()) .

The last inequality follows from the continuity of the trace operator : H1 () L2 () (recall

(3.6.1)).

Lemma 4.3.2. Let the assumptions on the data f, g, h stated at the beginning of Sect. 4.1 be

satisfied. In addition, suppose that the function ug be any lifting of the data g inside satisfying

kug kH1 () 2kgkH1/2 () (according to the definition (3.6.5) of the H1/2 ()-norm of g). Then,

one has

|F (v)| kf kL2 () + khkL2 (N ) kvkH1 () v H10,D () , (4.3.2)

|F (v)| kf kL2 () + khkL2 (N ) + 2C(A, a, a0 )kgkH1/2 () kvkH1 () v H10,D () . (4.3.3)

Proof. The first inequality follows immediately from the Cauchy-Schwarz inequality. Concerning

the second one, we have

Next, let us define a set of assumptions on the coefficients of the operator L, which ensure that

the bilinear form a(w, v) is coercive in H10,D () with respect to the H1 ()-norm. At first, let us

observe that

Z Z Z Z

2

a(v, v) = v Av v a v + a0 v + an v 2 v H10,D () . (4.3.4)

N

The first integral on the right-hand side is non-negative, since the assumption that the operator

is elliptic throughout the domain implies that A > 0 almost everywhere in , for any Rn .

However, this assumption is not sufficient to yield a control on the L2 -norm of v; we need a

stronger condition, expressed by the following definition.

Definition 4.3.3. The operator L is said to be uniformly elliptic in if there exists a constant

> 0 such that

Z Z

v Av kvk2Rn = kvk2(L2 ())n . (4.3.6)

The second integral on the right-hand side of (4.3.4) can be manipulated as follows. We first

note that vv = ( 12 v 2 ), so that

Z Z

v av = a ( 21 v 2 ) .

74 CHAPTER 4. ELLIPTIC PROBLEMS

Z Z Z

2 1

v av = 1

2a v + 2 an v 2 ;

the last integral is indeed an integral over N , since v vanishes on D . Thus, substituting this

expression and inequality (4.3.6) into (4.3.4) yields

Z Z

2

2 1

a(v, v) kvk(L2 ())n + 1

2 a + a0 v + 2 an v 2 v H10,D () .

N

At this point, we make the assumptions that a n 0 on N and that there exists a constant

such that 21 a + a0 almost everywhere in . Then, the above inequality implies

Finally, we observe that if the Poincare-Friedrichs inequality (3.7.4) holds for H10,D (), then it is

enough to assume 0 to get coercivity; indeed, recalling the first inequality in (3.7.6), we obtain

1+CP2 kvk2H1 () v H10,D () .

a(v, v) min(, ) kvk2(L2 ())n + kvk2L2 () = min(, ) kvk2H1 () v H10,D () .

Lemma 4.3.4. Assume that L be uniformly elliptic in , i.e., (4.3.5) holds true. Furthermore,

assume that a L () and that there exists a constant 0 for which 21 a + a0 almost

everywhere in ; let > 0 whenever the Poincare-Friedrichs inequality (3.7.4) does not hold for

H10,D (). Finally, assume that a n 0 on N . Then, the bilinear form a(w, v) is coercive in

H10,D (), with coercivity constant in the H1 ()-norm given by

2 if the Poincare-Friedrichs inequality holds in H10,D () ,

1+CP

= (4.3.7)

min(, ) if the inequality does not hold .

The three previous lemmas guarantee that the bilinear form a and the linear form F or F , which

define Problems 4.1.2 or 4.1.3, satisfy the assumptions of the Lax-Milgram Theorem. Consequently,

each of these Problem is well-posed. Recalling that the solution u of Problem 4.1.1 is given by

u = u0 + ug , hence in particular kukH1 () ku0 kH1 () + kug kH1 () , we arrive at the following final

result.

Theorem 4.3.5. Let the assumptions stated in Lemmas 4.3.1, 4.3.2 and 4.3.4 on the coefficients

A, a, a0 and the data f , g, h of the mixed Dirichlet/Neumann boundary-value problem (4.1.1) be

satisfied. Then, the weak formulation of the problem, given by Problem 4.1.1, admits one and only

one solution u, for which the following bound holds:

kukH1 () C kf kL2 () + kgkH1/2 () + khkL2 (N ) , (4.3.8)

where the constant C depends upon kAk(L ())nn , kak(L ())n , ka0 kL () and ka nkL (N ) .

4.3. WELL POSEDNESS OF ELLIPTIC PROBLEMS 75

i) The homogeneous Dirichlet problem for the Poisson equation, stated in (4.1.8), admits one

and only one variational solution u H10 () satisfying (4.1.9). Indeed, for Rthe operator one

has A = I, hence, = 1 in (4.3.5); consequently, the bilinear form a(w, v) = w v is coercive

in H10 (), with constant = 1/(1 + CP2 ), where CP is the Poincare-Friedrichs constant in H10 ()

(see (4.3.7)). Thus, from (4.2.10) and the fact that kf kH1 () kf kL2 () if f L2 (), we obtain

the following bound for the H1 ()-norm of u:

Note that a direct estimate on the norm kukH10 () = kuk(L2 ())n (see (recall Proposition 3.7.3)

can be obtained from the relations

which give

kukH10 () CP kf kL2 () . (4.3.10)

A similar well-posedness result holds if the data f is replaced by the more general data F

H1 (), according to (4.1.10). In this case, recalling (3.9.7), one has

n

!1/2

X

kukH1 () (1 + CP2 )kF kH1 () (1 + CP2 ) kfi k2L2 () . (4.3.11)

i=0

The latter result can be stated in the form of the following fundamental property.

: H10 () H1 ()

is an algebraic and topological isomorphism between these spaces.

ii) Consider now the non-homogeneous Neumann problem for the Helmholtz equation, stated

in (4.1.11). Assume that a0 almost everywhere in , for a suitable constant > 0. Then, the

problem admits one and only one variational solution u H1 () satisfying (4.1.12). The following

bound holds for the H1 ()-norm of u:

kukH1 () min(1, ) kf kL2 () + khkL2 () . (4.3.12)

iii) Finally, let us consider the mixed Dirichlet/Neumann problem for the convection-diffusion

equation, stated in (4.1.13). Let us assume that the Dirichlet part D of the boundary be non-

empty, and that the velocity field a be solenoidal, i.e., a = 0 in . Note that by the very

definition of D we have an 0 on N . Thus, the problem admits one and only one weak solution

u H1 () satisfying u = g on D and (4.1.14). The following bound holds for the H1 ()-norm of

u:

1+C 2

kukH1 () P kf kL2 () + kgkH1/2 () , (4.3.13)

76 CHAPTER 4. ELLIPTIC PROBLEMS

Theorem 4.3.5 ensures the existence and uniqueness of a weak solution of the boundary-value

problem (4.1.1), i.e., a solution of the weak formulation (4.1.5) of the problem. This formulation

has been derived by combining the partial differential equation to be satisfied inside the domain

with the Dirichlet and/or Neumann conditions to be satisfied at the boundary, via integrations

by parts.

From now on, we will move in the opposite direction, going back from (4.1.5) towards (4.1.1).

As a first step, we ask ourselves in which sense the weak solution u satisfies the partial differential

equation and the boundary conditions.

Let us begin with the equation. We recall that the most general framework to give sense to

a differential equation is the distributional sense. With this in mind, we observe that the space

H10,D () surely contains the space D() of all test functions for distributions. If we restrict (4.1.5)

to the functions v = in this space, we get

Z Z Z Z

(Au) u a + a0 u = f D() . (4.4.1)

This means that the equation is satisfied in the distributional sense, i.e., we have

From this relation, we can derive an additional information on the solution. Indeed, we write

(Au) = (au) + a0 u f in D () ,

and we observe that the right-hand side, which we write (a)u + a u + a0 u f , is a sum of

functions belonging to L2 (). Thus, the principal part of the operator, L(2) u = (Au), is

not just an element in D () (or in H1 ()), but is an element of L2 (). Therefore, if we define

the domain of the operator L(2) as the space

we conclude that the solution u is not just a function in H1 (), but it satisfies

u D(L(2) ) . (4.4.4)

The property (Au) L2 () (together with the property (au) L2 (), already used

above) allows us to write (4.4.1) in the equivalent form

Z Z Z Z

(Au) + (au) + a0 u = f D() ;

i.e., the equation is actually satisfied in a stronger sense than the distributional one, compare with

(4.4.2).

Condition (4.4.4) has also an important consequence on the interpretation of the Neumann

boundary conditions. Indeed, the following crucial property holds.

4.4. WHAT THE WEAK SOLUTION SATISFIES 77

v

Proposition 4.4.1. For any function v D(L(2) ), the conormal derivative n A

is well defined as

1/2 1/2

an element of the dual space H () = (H ()) . Precisely, the following formula holds:

Z Z

v

h , i = (Av) w + (Av) w H1/2 () , (4.4.6)

nA

Proof. (Sketch) The formula is the classical integration-by-parts formula if v and w are smooth

functions (e.g., if they belong to C1 ()). Then, the result is extended by a density argument.

We are ready to discuss in which sense the boundary conditions are satisfied. The Dirichlet

condition is clear: since u H1 (), its trace on is a function in H1/2 (), and we have required

that this function coincide with the data g on D . Concerning the Neumann condition, we first

observe that (4.4.4) and Proposition 4.4.1 yield

u

H1/2 () , (4.4.7)

nA

with Z Z

u

h ,v i = (Av) v + (Av) v v H10,D () ,

nA |

We now use (4.4.5) to write

Z Z

u

h ,v i = (( (au) + a0 u f ) v + (Av) v

nA |

Z Z

Z Z Z

= (Au) v u a v + a0 uv + an uv fv

N

Finally, we use

R the fact that u is the solution of the weak formulation (4.1.1), so that the right-hand

side equals N hv. In other words, the conormal derivative of u satisfies

Z

u

h , i = h H1/2 (), = 0 on N . (4.4.8)

nA N

conormal derivative of u induces a linear form on L2 (N ), which coincides with h. We write

u

=h in L2 (N ) ; (4.4.9)

nA N

this is precisely the way in which the Neumann boundary condition is satisfied by u.

We summarize the results obtained so far in the following theorem.

Theorem 4.4.2. Under the sole assumptions on the domain, the coefficients and the data for

which Theorem 4.3.5 holds, the weak solution of the boundary-value problem (4.1.1) is such that:

u u

H1/2 () and =h in L2 (N ) .

nA nA N

78 CHAPTER 4. ELLIPTIC PROBLEMS

At last, we investigate under which conditions on the domain , the coefficients A, a, a0 of

the operator and the data f, g, h the weak solution is a more regular function, up to being a

classical solution of the boundary-value problem (4.1.1).

We first discuss regularity inside the domain, then we deal with regularity up to the boundary.

The elliptic nature of the differential operator implies that there are no real characteristics, hence,

there is no propagation of singularities. In other words, the smootheness of u around a point in

only depends on the smoothness of the coefficients and the data f around that point. The result

is made precise in the following statement.

Proposition 4.5.1. Let be any open set contained with its closure in . Let m 0 be a

non-negative integer and assume that there exists an open set , satisfying , such that

Then,

u Hm+2 () .

The result means that if the coefficients of the operator are sufficiently smooth, then there is

a gain of two orders of Sobolev regularity between the data f and the solution u, i.e.,

f Hm ( ) u Hm+2 () . (4.5.1)

This is a manifestation of the property that elliptic operators are regularizing operators.

A partial yet simple justification of the previous result can be provided for the model equation

u = f in = Rn ,

using the powerful tool of the Fourier transform. Indeed, using (3.4.3), if both u and f belong to

L2 (Rn ) this equation is equivalent to

Z Z Z

1 + kk2(m+2) |u()|2 d = |u()|2 d + kk2m kk4 |u()|2 d

R n n n

ZR ZR

= 2

|u()| d + kk2m |f()|2 d

n n

ZR ZR

|u()|2 d + 1 + kk2m |f()|2 d .

Rn Rn

Using the expression (3.4.4) for the Sobolev norms, one easily gets

kukHm+2 (Rn ) C kukL2 (Rn ) + kf kHm (Rn ) , (4.5.3)

which is precisely (4.5.1) in the particular situation of being the full space.

4.5. BACK TO CLASSICS: THE REGULARITY OF THE WEAK SOLUTION 79

Going back to the general situation, condition u Hm+2 () with m large enough implies

classical regularity of u, thanks to the Sobolev Imbedding Theorem 3.8.2. Precisely, if m > n/2 2

then

u Hm+2 () u Ck, () ,

with k = [m + 2 n/2] and = m + 2 n/2 k if m n/2 is not an integer, or k = m + 1 n/2

and < 1 arbitrary if m n/2 is an integer. In particular, if the coefficients and the data f

are infinitely differentialble in (so that one can take m arbitrarily large), then u is infinitely

differentiable in .

It is worthwhile detailing certain results of minimal regularity in dimension two and three.

Corollary 4.5.2. Let n = 2 or 3, and let the assumptions of Proposition 4.5.1 hold. Then,

m = 0 implies u C0 (),

m = 1 implies u C1 (),

m = 2 implies u C2 ().

In the latter case, u is a classical solution of the partial differential equation at each point of .

Sobolev regularity up to the boundary of can be achieved if the coefficients and the data are

sufficiently smooth, and in addition if the boundary is a smooth manifold, which does not

contain points where a change between Dirichlet and Neumann boundary conditions occurs. The

simplest situation is described in the following proposition.

u Hm+2 () .

The result can be extended to more general situations, for instance when D and N are both

non-empty, but each connected component of is completely contained in one of these sets. In

all cases, a result of classical regularity similar to Corollary 4.5.2 holds in the whole of .

The question of assessing the regularity of the weak solution u up to the boundary becomes

quite delicate if is not a manifold of class Cm+1 (for instance, if it has corners), or if there is a

transition between Dirichlet and Neumann boundary conditions at some point of . We confine

ourselves to the illustration of some of the possible situations in the case of a polygonal domain.

80 CHAPTER 4. ELLIPTIC PROBLEMS

The following astonishingly simple example indicates how delicate is the issue of regularity in a

non-smooth domain. Consider the homogeneos Dirichlet problem for the Poisson equation, stated

in (4.1.8). Assume that is the square (0, 1) (0, 1), and take f to be the constant function

f = 1.

Despite the fact that we have a constant-coefficient operator and the data f and g = 0 are

2

constant, we cannot have u C2 () ! Indeed, if this were the case, we would have both xu2 and

2u

y 2

in C0 (). This would imply u C0 (), and since the laplacian of u equals 1 in , we

would also have u = 1 on ; in particular,

u(0, 0) = 1 .

2

On the other hand, the boundary condition on the side [0, 1] {0} implies xu2 (x, 0) = 0 for

2

0 < x < 1; by continuity, we would also have yu2 (0, 0) = 0. Similarly, using the boundary

2u

condition on the side {0} [0, 1] we would have x2

(0, 0) = 0, whence

u(0, 0) = 0 ,

We aim at getting some understanding of the effect of the presence of corners and/or transition

points between boundary conditions upon the regularity of the solution. We confine ourselves to

the following situation: we consider a simplified, yet significant form of the boundary-value problem

(4.1.1), i.e., the homogeneous mixed Dirichlet-Neumann problem for the Poisson equation,

u = f in ,

u = 0 on D , (4.5.4)

u = 0 on ,

N

n

and we pose the following question: under which conditions the implication

f L2 () = H0 () u H2 () (4.5.5)

is true?

We assume that is a bounded polygonal domain in R2 with vertices Vi , i = 1, . . . , I. Each

side i (which could even be a curved side) carries a unique type of boundary condition, i.e., either

i D or i N .

Let us assume that the vertex Vi is common to the sides i and i+1 (setting I+1 = 1 ); let

i (0, 2) be the measure of the angle at Vi , contained in . Note that the situation of a point

P internal to a side at which there is a change of type of boundary conditions can be included

into the present setting by considering P as an additional vertex of the polygon, with associated

angle of measure .

We also associate an angle i to each side i , by setting i = 0 if i N and i = 2 if

i D . (More generally, we could enforce as boundary condition the vanishing of an oblique

u

derivative of u, i.e., i

= 0 on i , where i is a fixed unitary vector which is not perpedicular to

the normal vector ni ; in such a case, i would be the measure of the angle between i and ni .)

With these notations at hand, we define the quantities

i i+1 m

i,m = , mZ, (4.5.6)

i

4.5. BACK TO CLASSICS: THE REGULARITY OF THE WEAK SOLUTION 81

and, correspondingly, the functions Si,m which, in a polar coordinate system (r, ) around the

vertex Vi , take the form

r i,m

Si,m (r, ) = cos(i,m i+1 ) i (r, ) ;

i i,m

here, each cut-off function i is infinitely differentiable, takes the value 1 at Vi and vanishes

identically if r is large enough.

The following result describes the structure of the solution of Problem (4.5.4).

Theorem 4.5.4. Let u H10,D () be the weak solution of Problem (4.5.4), with f L2 (). Then,

u can be represented as

u = ureg + using ,

where ureg H2 (), while

I

X X

using = ci,m (f ) Si,m

i=1 0<i,m <1

that 0 < i,m < 1; globally, u H2 () if all i,m fall outside the interval (0, 1).

Let us detail a few relevant situations.

Assume that both i and i+1 are contained in D , or in N . Then, i = i+1 , so that

i,m = m .

i

If is convex around Vi , i.e., if 0 < i < , then i > 1, so that there is no m Z such that

0 < i,m < 1: the solution is H2 in a neighborhood of Vi .

This result is a particular case of the following property.

L2 ().

Obviously, the counter-example at the beginning of the present section shows that a similar

result cannot hold for higher-order regularity: the convexity of the domain is not enough to ensure

that f H2 () implies u H4 ().

Going back to our discussion, if is not convex around Vi , i.e., if < i < 2, then 12 < i < 1,

so that there is exactly one m Z, namely m = 1, such that 0 < i,m < 1: the solution is not

H2 in any neighborhood of Vi . For instance, if i = 23 as in the re-entrant corner of an L-shaped

domain, then one can prove that u belongs at most to H7/4 in a neighborhood of Vi .

ii) Consecutive sides carrying different boundary conditions

Assume that i D , whereas i+1 N . Then, i i+1 = 2 , so that

1

i,m = m .

2 i

82 CHAPTER 4. ELLIPTIC PROBLEMS

switches from Dirichlet to Neumann, then i = , so that i,m = 12 m. Thus, 0 < i,m < 1

exactly for m = 0, and we have a singularity at Vi : one can prove that u belongs at most to H3/2

in a neighborhood of Vi .

Conversely, if i and i+1 are perpendicular to each other, i.e., if i = 2 , then i,m = 1 2m.

In this case, no i,m satisfies 0 < i,m < 1, hence, the solution is H2 in a neighborhood of Vi .

4.6 Exercises

4.1. Consider the following problem in = (0, 1)2 IR2 :

u = f in

u = 0 on 2 3 4

u + u = 0 on 1

n

where IR and 1 = (0, 1) {0}, 2 = {1} (0, 1), 3 = (0, 1) {1}, 4 = {0} (0, 1).

(ii) Find the conditions on which guarantee the coercivity of the associated bilinear form.

4.2. Setting up a suitable bilinear form in V = H10 () H10 (), use the Lax-Milgram Theorem to

prove the existence and uniqueness of the solution of the following elliptic system:

u2

u1 + u1 + = f1 in

x1

u2 + u1 + u2 = f2 in

u1 = u2 = 0 on

Chapter 5

Under the name of Maximum Principle, we find several important theoretical properties of the

solutions of elliptic and parabolic problems, all related to the ordering relation between real

numbers. The Maximum Principle can be expressed in various forms, from the classical ones to

the more general statements derived from the weak, or variational, formulations of the problems.

In this chapter, we will confine ourselves to elliptic boundary value problems. The Maximum

Principle for parabolic problems will be briefly accounted for in Sect.

In this section, for the sake of simplicity we consider the Laplace operator only, although most of

the results hold for more general elliptic operators, under suitable assumptions on the coefficients.

We general treatment is postponed to the next section.

The first expression of the Maximum Principle that we present has an immediate physical

interpretation. Let a thin elastic membrane occupy, when no load is applied to it, the position of

a plane domain , and let it be attached to the boundary . Then, if we apply a vertical load

pointing upwards, the membrane gets inflated upwards as well.

Proposition 5.1.1. Let be a bounded domain, with sufficiently smooth boundary . Let

u C0 () C2 () be the solution of the problem

u = f in

u = 0 on .

If f 0 in , then u 0 in .

Proof. Suppose at first f > 0 in and assume by contradiction that there exists x such

that u(x ) < 0. Since u is continuous in the bounded domain , there exists x such that

u(x ) = min u(x ) < 0.

x

Then, necessarily u(x ) = 0, and furthermore the Hessian of u in x is nonnegative, which implies

in particular

2u

0 i = 1, 2, . . . , n.

x2i

Thus

n

X 2u

u(x ) = (x ) 0

i=1

x2i

83

84 CHAPTER 5. THE MAXIMUM PRINCIPLE

Suppose now f 0 in ; take any > 0 and set f (x ) = f (x ) + , so that f > 0 in . Let

u be the solution of the problem

u = f in ,

u = 0 on .

By linearity we can use the superposition principle, which allows us to write

u (x ) = u(x ) + u(x ), (5.1.1)

where u solves the problem

u = 1 in ,

u = 0 on .

This solution satisfies u C0 () C2 (), since the right-hand side is infinitely smooth and is

supposed to be sufficiently smooth to allow this. Consequently, by (5.1.1) the same property holds

for u . Thus, we can apply to this function the result proven above and get u 0 in . Taking

the limit in (5.1.1) as 0, we finally get u 0 in .

Remark 5.1.2. The proof of the Proposition easily shows that u cannot have points of strict

local minimum inside the domain.

Remark 5.1.3. The previous result is also known as the Weak Maximum Principle. It is possible

to prove even a Strong

R Maximum Principle, which states that if is a connected domain and if

f 0 in with f (x ) dx > 0, then u > 0 in .

So far, we have assumed zero boundary conditions and a non-zero right-hand side. The fol-

lowing result applies to harmonic functions, and provides a bound for their values at the interior

of the domain in terms of the boundary values.

Proposition 5.1.4. Let be a bounded domain and let u C0 () C2 () be the solution of the

problem

u = 0 in ,

u = g on ;

then

min g u(x) max g, x .

Proof. Suppose that the second inequality does not hold, and assume that there exists x

such that

u(x ) = max u > max g;

we can then choose IR such that

max g < < u(x ).

Define now := {x : u(x ) > }; this set is a nonempty (since x ), bounded (since

) and open (since u is continuous); furthermore, = {x : u(x ) = }. In , u

solves the Dirichlet problem

u = 0 in ,

u = on .

By the uniqueness of the solution, we must then have u in ; but, by assumption, it results

u(x ) > , a contradiction. The first inequality in the thesis is proven similarly.

5.2. VARIATIONAL RESULTS 85

Remark 5.1.5. For harmonic functions, we can actually state more: a harmonic function cannot

have strict maxima or minima inside its domain.

The result is an immediate consequence of the following property of harmonic functions, which

we will derive in a moment: given a point x and any neighborhood BR (x ) = {z : kz x k < R}

contained in , with boundary R = {z : kz x k = R}, we have the expression:

Z

1

u(x ) = u() d, (5.1.2)

|R | R

where |R | denotes the measure of R . As a consequence, u cannot achieve, for instance, a local

maximum value in x , since in that case we would have u() < u(x ) for every which parametrizes

R , and then Z Z

1 1

u(x ) = u() d < u(x ) d = u(x ),

|R | R |R | R

a contradiction. For a local minimum value in x , the reasoning is analogous.

We will derive (5.1.2) in dimension 2, the general case being similar. Consider the function

1 r

v(z ) = log , with r = kz x k, which according to (2.3.2) satisfies

2 R

v = x in BR (x ),

v = 0 on R .

Then,

Z Z

v

u(x ) = hx , ui = hv, ui = v u dx u d

BR (x ) R n

Z Z Z Z

v u v

= vu dx u d + v d = u d,

BR (x ) R n R n R n

v 1

since u is harmonic in BR (x ) and v vanishes on R . We conclude by observing that =

n 2R

on R .

The results presented in the previous section are special cases of a very general theorem, which can

be proved by a variational technique due to G. Stampacchia. In order to obtain such a theorem,

we need at first some preliminary considerations.

Let v be a function defined in ; we set

The functions v + and v are said to be the positive part and the negative part of v, respectively.

Note that the following decomposition holds true:

v = v+ v ;

86 CHAPTER 5. THE MAXIMUM PRINCIPLE

measure; for them, the pointwise value is meaningless. Therefore, an expression of the form v w

in should be understood as v w almost everywhere in .

Let us now consider a generic second order elliptic operator, say

Lu = (Au) + a u + a0 u,

Z Z Z

a(u, v) = Auv dx + a u v dx + a0 uv dx (5.2.1)

is continuous and coercive in H10 (); furthermore, given f L2 () and g H1/2 () C0 (),

let us denote by u the solution of the Dirichlet problem

Lu = f in ,

u = g on ,

that is (

u H1 (), u| = g,

a(u, v) = (f, v) v H10 ().

Finally, let us set

mg = min g and Mg = max g. (5.2.2)

i) if f a0 mg 0 in , then u mg in ;

ii) if f a0 Mg 0 in , then u Mg in .

Proof. We prove i), since the proof of ii) is similar. Let us write u in the form u = (u mg ) + mg

and let us substitute it in the variational formulation:

a(u mg , v) = (f a0 mg , v).

and furthermore it results (u mg )| = g mg 0 on , hence, (u mg )

| = 0; therefore

(u mg ) H10 (),

a (u mg )+ , (u mg ) a (u mg ) , (u mg ) = f a0 mg , (u mg ) . (5.2.4)

5.2. VARIATIONAL RESULTS 87

a v+ , v = 0

since the supports of v + and v do not intersect in , except for possible sets of measure zero;

consequently, equation (5.2.4) becomes

a (u mg ) , (u mg ) = f a0 mg , (u mg ) .

k(u mg ) k21, f a0 mg , (u mg ) ;

f a0 mg , (u mg ) 0,

which implies

k(u mg ) k21, 0

and finally (u mg ) = 0. This means that u mg 0 in , i.e., u mg in .

Remark 5.2.3. Instead of g C0 (), one could make the weaker assumption that g L ().

In that case, the theorem holds after replacing (5.2.2) by

Remark 5.2.4. By inspecting the proof, it is easily seen that the implications i) and ii) of the

theorem hold if mg indicates any number min g and Mg indicates any number max g.

In the examples below we shall show how Stampacchias Theorem can be applied to some

particular cases of elliptic problems.

Example 5.2.5. Let us consider the Dirichlet problem for the Laplace operator

u = f in ,

u = 0 on ;

proven in the classical case (cf. Proposition 5.1.1).

Instead, if u = g 6= 0 on , then we get u min g in . k

u = 0 in ,

u = g on ,

is such that a0 = 0 and f 0; from Stampacchias Theorem it follows that

(i) f 0 u ming

min g u max g,

(ii) f 0 u maxg

again a result already proved in the classical case (cf. Proposition 5.1.4). k

88 CHAPTER 5. THE MAXIMUM PRINCIPLE

Example 5.2.7. Let us consider the Dirichlet problem for the Helmholtz operator

u + u = f in ,

u = g on .

We have a0 = 1. Assume at first that f 0 in and mg 0: we conclude that u min g in

. On the other hand, let again f 0 in but now assume mg > 0; then, the assumptions of

Stampacchias Theorem may not be satisfied and we can indeed have u < min g in , as shown

in the forthcoming discussion of singular perturbation problems. However, in this case we can

apply Remark 5.2.4 with mg = 0 and get at least u 0 in . k

The hypothesis of Stampacchia Theorem involve neither second-order nor first-order coefficients

of the operator L. Therefore, these coefficients can be very small or large (provided the coercivity

of the form a is guaranteed) and the result still holds.

As an example, let us consider the so-called singular perturbation problem

u + u = f in ,

u = g on ,

where > 0 is a constant. This is an extremely simplified example of a reaction-diffusion problem,

in which the second-order part of the operator models a diffusion process (say, by mechanical

molecular interactions), whereas the zeroth-order part models a reaction process (say, by the

action of chemical agents). If the diffusion coefficient is small, then the solution u = u , far from

the boundary, is close to the function f , as if the equation were simply u = f without the term

u; u eventually adjusts itself to the boundary condition, i.e., the second-order term becomes

influential, only in a small transition region attached to the boundary. This region, called the

boundary layer, has length of order in the direction normal to the boundary.

To illustrate this behavior, let us consider the one-dimensional problem

u + u = f in (0, 1),

u(0) = g0 ,

u(1) = g1 ,

with f, g0 , g1 IR. In order to get the solution, let us observe that the function w = u f solves

the problem

w + w = 0 in (0, 1),

w(0) = g0 f,

w(1) = g1 f.

By linearity and the superposition principle, w can be expressed as w = (g0 f )u0 + (g1 f )u1 ,

where u0 is the solution of the problem

u0 + u0 = 0 in (0, 1),

u0 (0) = 1,

u0 (1) = 0,

whereas u1 is the solution of the problem

u1 + u1 = 0 in (0, 1),

u1 (0) = 0,

u1 (1) = 1.

5.2. VARIATIONAL RESULTS 89

In conclusion, we have

u(x) = f + (g0 f )u0 (x) + (g1 f )u1 (x).

e(1x)/ e(1x)/ ex/ ex/

u0 (x) = and u1 (x) = 1/ ,

e1/ e1/ e e1/

from which it is easily seen that u0 is very close to 0 for c < x < 1, whereas u1 is very close to

0 for 0 < x < 1 c , where c > 1 is any large enough constant independent of .

On the whole, the solution u approximately equals f in (c , 1 c ), and exhibits two

boundary layers in (0, c ) and (1 c , 1) (see Fig. ....)

Going back to the maximum principle discussed at the end of the previous section, we note

that if 0 < f < min(g0 , g1 ), then we will also have 0 < u(x) < min(g0 , g1 ) at some point x (0, 1),

provided is small enough.

provided by the advection-diffusion problem

u + a u = f in ,

u = g on ,

where > 0 is again a constant and a is a smooth vector field such that kak 1 in . The term

a u models the transport of the scalar quantity u (which may represent the temperature of a

fluid, or the concentration of a pollutant in a fluid) along the streamlines of the field a.

If is small, so that the term u may be neglected, the equation in reduces to a particular

case of the linear first order equation (1.3.1) considered in Sect. 1.3. We have seen that this

equation may be solved with the boundary condition assigned on the inflow boundary defined

as in (1.3.4). The solution u = u , far from the portion of the boundary 0 + , is close to

the solution u of the reduced first-order problem

a u = f in ,

u = g on .

The transition occurs in boundary layers of differents size: the length in the direction normal

to the boundary is of order for the boundary layer attached to the characteristic boundary

0 , whereas it is of order (thus, much smaller) for the boundary layer attached to the outflow

boundary + .

Again we use a one-dimensional problem to illustrate the structure of the solution. Let us

consider

u + u = f in (0, 1),

u(0) = g0 ,

u(1) = g1 ,

u = f in (0, 1),

u(0) = g0 ,

90 CHAPTER 5. THE MAXIMUM PRINCIPLE

(note that x = 0 is the inflow boundary point), getting u(x) = g0 + f x. Next we make the change

of variable w = u u and we observe that w is the solution of

w + w = 0 in (0, 1),

w(0) = 0,

w(1) = g1 g0 f.

By linearity we get

u(x) = g0 + f x + (g1 g0 f )u1 (x),

where now u1 is the solution of

u1 + u1 = 0 in (0, 1),

u1 (0) = 0,

u1 (1) = 1.

ex/ 1 e(x1)/ e1/

u1 (x) = = ,

e1/ 1 1 e1/

which shows that u1 is very close to 0 for 0 < x < 1 c, for any c > 1 large enough.

On the whole, the solution u approximately equals the linear function g0 + f x in (0, 1 c)

and exhibits an outflow boundary layer in (1 c, 1) (see Fig. ....)

Concerning the maximum principle, we have a0 = 0, hence, Stampacchias Theorem yields

u min(g0 , g1 ) in if f 0, and u max(g0 , g1 ) in if f 0. Note that we have min(g0 , g1 )

u max(g0 , g1 ) in if f = 0.

Chapter 6

Self-adjoint Problems

6.1 Introduction

In this chapter, we introduce the concept of eigenvalue and eigenfunction of a uniformly elliptic,

self-adjoint operator with appropriate boundary conditions. They can be considered as the gener-

alization to an infinite dimensional Hilbert space H of the concept of eigenvalue and eigenvector of

a symmetric positive-definite matrix in a Euclidean space Rn . The eigenfunctions of the operator

form an orthonormal basis in H, with respect to which the operator is diagonalized.

as follows. Consider the wave equation

2u

u = f

t2

in a bounded domain Rn , submitted to honogeneous Dirichlet boundary conditions

u=0.

In one dimension (n = 1), u can be interpreted as the (small) displacement of a guitar string in the

direction perpendicular to the plane of the guitar, under the effect of a density force; the boundary

condition forces the string to remain attached to the guitar. The two-dimensional analog is the

displacement of a drum membrane (or drum skin) in the direction perpendicular to the drum

surface. (For simplicity, we have set to 1 all physical constants.)

An important characteristic of the musical instrument is represented by the set of the modes

of free vibration of the guitar string or the drum skin. They correspond to those solutions of the

wave equation with no forcing term (f = 0), which are periodic in time. More precisely, they can

be defined as the solutions of the form

where R is the frequency of pulsation, whereas w determines the spatial shape of the pulsation.

Differentiating and substituting into the wave equation, we get

eit 2 w(x) w(x) = 0 ,

91

92 CHAPTER 6. SPECTRAL THEORY FOR ELLIPTIC SELF-ADJOINT PROBLEMS

w = 2 w in ;

on the other hand, the boundary condition for u is equivalent to the boundary condition for w,

w=0 on .

Problem 6.1.1. Find all real numbers and all non-zero functions w defined in such that

(

w = w in ,

(6.1.1)

w = 0 on ,

Using the functional machinery developed in Chap. 4, this problem can be equivalently written

in variational form as follows.

Problem 6.1.2. Find all R and w H10 (), w 6= 0, such that

Z Z

w v = w v v H10 () . (6.1.2)

The situation described so far is just an example, and it can be generalized in various ways,

e.g. by considering variable-coefficient elliptic operators or other types of boundary conditions. In

the next section, we will provide an abstract framework to study spectral problems formulated in

a variational form.

Consider a Gelfand triple (V, H, V ) of separable Hilbert spaces (see (3.9.1)), and assume in addi-

tion that the inclusion V H is compact, i.e., from every sequence {vn }n1 which is bounded in

V , one can extract a subsequence {vnj }j1 which is convergent in H.

Denote by (u, v) the inner product in H, and let a : V V R be a bilinear form which is

continuous and coercive in V , with coercivity constant > 0 (i.e., it satisfies the assumptions of

the Lax-Milgram theorem); assume in addition that the form a is symmetric.

Let us consider the following eigenvalue problem.

Problem 6.2.1. Find all R and w V , w 6= 0, such that

The solutions of this problem are referred to as the eigenvalues and the eigenfunctions w of

the bilinear form a. Note that any eigenfunction is defined up to a multiplicative constant, i.e., if

w is an eigenfunction, then w is also an eigenfunction, for any real 6= 0. Furthermore, given

an eigenfunction w, the corresponding eigenvalue can be expressed as the Rayleigh quotient

a(w, w)

= .

(w, w)

Problem 6.2.1 can be solved by resorting to the spectral theory for a compact, self-adjoint and

positive operator in a Hilbert space. Let us recall the main result of this theory.

6.2. THE ABSTRACT VARIATIONAL EIGENVALUE PROBLEM 93

Theorem 6.2.2. Let X be a separable Hilbert space, with inner product (x, y). Let T : X X

be a linear operator, satisfying the following assumptions:

a) T is compact, i.e., from every bounded sequence {xn }n1 , one can extract a subsequence

{xnj }j1 such that {T xnj }j1 is convergent;

Under these assumptions, there exist a sequence {k }k1 of strictly positive real numbers and a

sequence {xk }k1 of elements of X, which satisfy the relations

T xk = k xk k = 1, 2, . . . (6.2.2)

function). Eigenvalues and eigen-elements enjoy the following properties:

i) The sequence of eigenvalues is non-increasing and converging to 0, i.e.,

k

(xk , x ) = k , k, 1 . (6.2.4)

iii) This system is indeed an orthonormal basis in X, i.e., every x X can be uniquely repre-

sented as

X

x= k xk with k = (x, xk ) ; (6.2.5)

k=1

the coefficients k are termed the (generalized) Fourier coefficients of x with respect to the

basis of eigen-elements of T . The convergence of the series has to be meant in the norm of

X, i.e.,

XN

kx xN k 0 as N , where xN = k wk .

k=1

iv) The following representation of the norm in X, termed Parseval identity, holds:

X

2

kxk = |k |2 x X . (6.2.6)

k=1

V () = {z X : T z = z}

is a vector space, termed the eigenspace of . As a consequence of (6.2.3), V () has finite dimen-

sion, say m 1 if we have for some k

m.

94 CHAPTER 6. SPECTRAL THEORY FOR ELLIPTIC SELF-ADJOINT PROBLEMS

Remark 6.2.4. Recall that for a linear operator, the property of compactness is stronger than

the property of continuity, i.e., if T is compact, then necessarily it is continuous. Indeed, by

contradiction, if there were no constant C > 0 such that

kT xk C kxk x X ,

kT xn k > n kxn k ;

Clearly, no converging subsequence can be extracted from {T yn }, a contradiction with the property

of compactness.

We are ready to introduce the operator T : H H associated with our Problem 6.2.1.

Proposition 6.2.6. The operator T defined above satisfies assumptions a)-c) of Theorem 6.2.2.

Proof. Let us check compactness. Let C > 0 denote the continuity constant of the inclusion

V H, i.e., kvkH CkvkV for all v V . Furthermore, let F V be the form associated with

f , i.e., the form defined by F (v) = (f, v) for all v V . Then, recalling (4.2.10) and the definition

of dual norm (see Sect. 3.9), one has

1 1

kukH kukV kF kV Ckf kH .

C

2

This implies both kT f kH C kf kH (i.e., the continuity of T , which we have seen above to be a

necessary condition for compactness), and

C

kT f kV kf kH .

Thus, if {fn } is a bounded sequence in H, then {T fn } is a bounded sequence in V and since the

inclusion V H is compact by assumption, we can extract a subsequence {T fnj } converging in

H. This proves that T is compact.

In order to prove assumption b), consider any f and g in H, and set u = T f , w = T g. Using

the symmetry of the inner product in H and the definition of T , we have

Finally, the positivity of T immediately follows from the coercivity of a. Indeed, for any f H,

one has

(T f, f ) = (f, T f ) = (f, u) = a(u, u) kuk2V 0 ,

6.2. THE ABSTRACT VARIATIONAL EIGENVALUE PROBLEM 95

H.

The result just proven ensures us that statements i)-iv) of Theorem 6.2.2 apply to our operator

T . Let us rephrase them in the language of the variational setting.

We will denote by wk , k = 1, 2, . . . , the eigenfunctions of T . Then, each T wk satisfies

a(T wk , v) = (wk , v) v V ;

yields

k a(wk , v) = (wk , v) v V .

Setting

1

k = ,

k

we obtain

a(wk , v) = k (wk , v) v V . (6.2.8)

We arrive at the conclusion that T wk = k wk holds if and only if k R+ and wk V are a

solution of the variational eigenvalue problem (6.2.1). We collect the results obtained so far in the

following theorem.

Theorem 6.2.7. Under the assumptions on the Gelfand triple (V, H, V ) and the bilinear form

a(u, v) stated at the beginning of this section, the variational eigenvalue Problem 6.2.1 admits a

sequence of real strictly positive eigenvalues k and corresponding eigenfunctions wk V , k =

1, 2, . . . , with the following properties:

k

(wk , w ) = k , k, 1 . (6.2.10)

iii) This system is indeed an orthonormal basis in H, i.e., every v H can be uniquely repre-

sented as

X

v= vk wk with vk = (v, wk ) ; (6.2.11)

k=1

X

kvk2H = |vk |2 v H . (6.2.12)

k=1

Corollary 6.2.8. The previous conclusions apply to the eigenvalue Problem 6.1.2 for the operator

with Dirichlet boundary conditions in a bounded domain .

96 CHAPTER 6. SPECTRAL THEORY FOR ELLIPTIC SELF-ADJOINT PROBLEMS

Proof. We know that (H10 (), L2 (), H1 ()) is a Gelfand triple (see Sect. 3.9), and that the

inclusion H10 () L2 () is compact if is bounded (by Rellich Theorem 3.8.1).

We close this section by showing that the system of eigenfunctions {wk }k1 provides a basis

not only in H but also in other functional spaces related to our problem, such as V and V ; in

addition, the norm of an element in such a space can be expressed in terms of a suitably weighted

norm of the sequence of its (generalized) Fourier coefficients. Let us begin with the space V .

Setting v = w in (6.2.8) and using (6.2.10), we obtain

a(wk , w ) = k k , k, 1 . (6.2.13)

This means that the eigenfunctions form an orthogonal system in V , for the inner product (u, v)a =

a(u, v) associated with the bilinear form a (which induces a norm, kvka , uniformly equivalent to

the norm kvkV , see Remark 4.2.7). The system is indeed a basis for V , as stated by the following

result.

X

vV if and only if k |vk |2 < + .

k=1

In this case, the series (6.2.11) converges also in V and the following Parseval identity holds true:

X

a(v, v) = kvk2a = k |vk |2 . (6.2.14)

k=1

Proof. For any fixed N 1, let VN = span {wk : 1 k N } the subspace of V spanned by the

first N eigenfunctions.

PN that v V . Define vN VN as the truncation to N terms of the series (6.2.11), i.e.,

Assume

vN = k=1 vk wk . It is easily seen, thanks to (6.2.13), that vN is the orthogonal projection of v

upon VN in the inner product (u, v)a , i.e., it satisfies

(v vN , z)a = 0 z VN .

Pythagoras theorem in the norm k . ka applies, yielding

which implies

N

X

k |vk |2 = kvN k2a kvk2a < + , N 1 .

k=1

Thus, the series in the statement of the proposition is convergent, and the other properties follow

immediately from this result.P PN

Conversely, assume that 2

k=1 k |vk | < +. Then, setting vN = k=1 vk wk for all N 1,

we have that the sequence {vN } is a Cauchy sequence in V , since

N

X

kvN vM k2a = k |vk |2 0 if N >M .

k=M +1

6.2. THE ABSTRACT VARIATIONAL EIGENVALUE PROBLEM 97

by the previous theorem, necessarily v = v by the uniqueness of the limit. We conclude that

v V.

Proposition 6.2.9 provides an example of the fundamental property that for an element v H

the condition of being more regular (in a suitable sense) is equivalent to a faster decay of its

Fourier coefficients vk as k , since the sequence {k } is diverging to + (recall (6.2.9)). Note

in particular that from v belonging to H we can only infer that

since the series (6.2.12) is convergent; on the other hand, from v belonging to V we infer that

1

vk = o k.

k

Another example of this property is as follows. Recall the definition (4.2.6) of the operator

A : V V associated with the bilinear form a. Note that the relations (6.2.8) can be equivalently

written as

Awk = k wk , k = 1, 2, . . . .

Let us introduce the domain of the operator A as the subspace of V

D(A) = {v V : Av H} . (6.2.15)

P

One can prove (see the remark below) that if v =

k=1 vk wk belongs to D(A), then

X

X

Av = vk Awk = k vk wk ,

k=1 k=1

so that

X

kAvk2H = 2k |vk |2 < + .

k=1

In this case,

1

vk = o k.

2k

In the model situation of the Laplacian with Dirichlet boundary conditions, we have H =

L (), V = H10 () and D() = H10 () H2 () provided has a C1 boundary (Proposition 4.5.3)

2

or is convex (Property 4.5.6). In one of these cases, the following equivalences hold:

X

v L2 () if and only if |vk |2 < + , (6.2.16)

k=1

X

v H10 () if and only if k |vk |2 < + , (6.2.17)

k=1

X

v H10 () H2 () if and only if 2k |vk |2 < + . (6.2.18)

k=1

On can continue this sequence indefinitely, by considering domains of the successive powers m

of the Laplacian.

98 CHAPTER 6. SPECTRAL THEORY FOR ELLIPTIC SELF-ADJOINT PROBLEMS

Remark 6.2.10. We can also expand in eigenfunction series any element F of the dual space V .

Indeed, we define its (generalized) Fourier coefficients by setting

Fk = F (wk ) , k = 1, 2, . . . ;

P

then, for any v = k=1 vk wk V , one has (at least formally)

X

X

F (v) = vk F (wk ) = Fk vk . (6.2.19)

k=1 k=1

Since F (v) is finite for all v V , one can prove using Proposition 6.2.9 that necessarily

X 1

|Fk |2 < + ,

k

k=1

and actually the square root of this expression is the dual norm of F , if V is equipped with the

norm kvka . In other words, one can prove that

X

X 1

F = Fk wk V if and only if |Fk |2 < + ;

k

k=1 k=1

in this case, the series of F is convergent in V , and the action of F on any v V can be expressed

as in (6.2.19).

We have seen (Corollary 6.2.8) that in a bounded domain Problem 6.1.2 admits a sequence {k }k1

of strictly positive eigenvalues, diverging to + as k , with corresponding eigenfunctions

{wk }k1 , which form an orthonormal basis in L2 ().

The following result is of some interest, since it relates the smallest eigenvalue 1 to the

Poincare constant CP () of the domain (defined in Proposition 3.7.2).

Property 6.3.1. The Poincare constant of the domain is given by

1

CP () = .

1

Proof. We recall that the Poincare constant of the domain is the smallest constant CP for which

(3.7.1) holds; this condition can be written as

R R

with (u, v) = uv and a(u, v) = u v. Recalling (6.2.12) and (6.2.14), we obtain

X

X

2

|vk | CP2 k |vk |2 ;

k=1 k=1

X

X

2 k

|vk | CP2 1 |vk |2 .

1

k=1 k=1

6.3. CLASSICAL EXAMPLES. SEPARATION OF VARIABLES 99

This inequality is true if CP2 1 = 1, since k /1 1 for any k by (6.2.9). On the other hand, with

such a choice the inequality becomes an equality for v = w1 , meaning precisely that

1

CP =

1

is the smallest admissible constant in (3.7.1).

Next, we explicitly compute the eigenvalues and eigenfunctions of Problem 6.1.2 in some rele-

vant cases.

We consider the one-dimensional version of Problem 6.1.1, i.e., we look for all the solutions of the

equations (

w = w in (0, L) ,

(6.3.1)

w(0) = w(L) = 0 ,

where L > 0 is fixed. The equation in (0, L) is a second-order, linear, constant-coefficient ordinary

differential equation, which is well known to admit the general integral

w(x) = Aei x

+ Beix

,

w(x) = a sin( x) + b cos( x) , a, b R arbitrary .

Condition w(0) = 0 forces b = 0, whereas condition w(L) = 0 yields L = k for any integer

k > 0 (since the left-hand side is strictly positive). Thus, we have found the eigenvalues of our

problem:

2

k = k 2 2 , k = 1, 2, . . . , (6.3.2)

L

with corresponding eigenfunctions

wk (x) = ak sin(k x) .

L

The parameter ak is determined by the normalization condition

Z L Z L

2 2

wk (x) dx = ak sin2 (k x) dx = 1 ,

0 0 L

p

which easily gives ak = 2/L. Thus, the eigenfunctions of our problem are

r

2

wk (x) = sin(k x) , k = 1, 2, . . . , (6.3.3)

L L

Following the same procedure illustrated above, one can easily compute (see Exercise 6.1)

eigenvalues and eigenfunctions of other boundary-value problems for the second derivative opera-

tor, such as a mixed Dirichlet/Neumann boundary-value problem

w = w in (0, L), w = w in (0, L),

w(0) = 0, or w (0) = 0, (6.3.4)

w (L) = 0, w(L) = 0.

100 CHAPTER 6. SPECTRAL THEORY FOR ELLIPTIC SELF-ADJOINT PROBLEMS

(

w = w in (0, L) ,

(6.3.5)

w (0) = w (L) =0,

Remark 6.3.2. Note that one cannot directly apply Theorem 6.2.7 to the variational formulation

of the pure Neumann problem (6.3.5), i.e.,

Z L Z L

w v dx = w v dx v V = H1 (0, L) , (6.3.6)

0 0

RL

since the bilinear form a(w, v) = 0 w v dx is not coercive in V : one has a(w, w) = 0 for all

constant functions w. In other words, the problem admits the eigenvalue = 0 (with eigenfunction

w = const), which is excluded by Theorem 6.2.7.

In order to apply the theorem, one uses the trick of adding the L2 -inner product on both sides

of (6.3.6), i.e., one considers the modified problem

Z L Z L Z L

w v dx + w v dx = w v dx v V = H1 (0, L) ,

0 0 0

with = + 1. Now the bilinear form on the left-hand side is precisely the inner product in V ,

so that this problem fulfils the assumptions of Theorem 6.2.7. The eigenfunctions wk of the two

problems are the same, whereas the eigenvalues of (6.3.6) are given by k = k 1.

The same trick of shifting the eigenvalues applies whenever one has to solve a general eigenvalue

Problem 6.2.1, in which the bilinear form a(w, v) is not coercive in V , but is such that the shifted

form a(w, v) + (w, v) is indeed coercive for > 0 large enough.

Next, we solve Problem 6.1.1 in the square = (0, L)2 . The nature of the domain (a cartesian

product of intervals) and the form of the equation (a constant-coefficient operator) suggests to

adopt the separation of variables approach, which consists of looking for a solution in the form

i.e., a product of a function of the x-variable alone and a function of the y-variable alone. Substi-

tuting into the differential equation w = w, we get

(x) (y)

=. (6.3.8)

(x) (y)

This identity holds in , i.e., for all x (0, L) and (independently) for all y (0, L). Keeping y

fixed and varying x, we see that necessarily

(x)

= constant (say, ) ,

(x)

6.3. CLASSICAL EXAMPLES. SEPARATION OF VARIABLES 101

(y)

= constant (say, ) .

(y)

We deduce that (6.3.8) is equivalent to

(x) (y)

=, =, =+ . (6.3.9)

(x) (y)

Enforcing the boundary condition w = 0 on the vertical side x = 0 of the square yields

(0)(y) = 0 , 0yL,

and since cannot be identically 0 (otherwise w would be so), then necessarily (0) = 0. Con-

sidering all other sides of the square, we find by similar arguments (L) = 0, (0) = 0 and

(L) = 0.

In conclusion, a function w of the form (6.3.7) is an eigenfunction of our problem if and only

if and satisfy

( (

= in (0, L) , = in (0, L) ,

(0) = (L) = 0 , (0) = (L) = 0 ,

i.e., if and only if both and are arbitrary eigenfunctions of the one-dimensional problem

considered in the previous subsection. In this way, we find the eigenfunctions

2

whk (x, y) = sin(h x) sin(k y) , h, k = 1, 2, . . . , (6.3.10)

L L L

with associated eigenvalues

2

hk = (h2 + k2 ) , h, k = 1, 2, . . . , (6.3.11)

L2

A reasonable question at this point is whether there exist eigenfunctions other than those found

so far. The answer is no, since one can prove that the set (6.3.10) is complete in L2 (), so that

any w orthogonal to all whk would necessarily be the zero function in L2 ().

Note that the adopted labeling of eigenfunctions and eigenvalues, by two indices h and k,

differs from the single-index labeling used in Theorem 6.2.7. Yet, it is not difficult to see that the

set {(h, k) : h, k = 1, 2, . . . } can be numbered in such a way that the monotonicity conditions

(6.2.9) are fulfilled.

At last, we solve Problem 6.1.1 in the unit circle centered at the origin, i.e., in = B(0, 1). It is

natural to resort to polar coordinates (r, ); recalling (2.3.1), the differential equation w = w

becomes 2

w 1 w 1 2w

+ + = w .

r 2 r r r 2 2

The domain is transformed into the product of intervals [0, 1) [0, 2), which suggests as above

to look for solutions in separated form

102 CHAPTER 6. SPECTRAL THEORY FOR ELLIPTIC SELF-ADJOINT PROBLEMS

1 1

(r)() + (r)() + 2 (r) () = (r)() ;

r r

dividing by (r)(), we get

(r) 1 (r) 1 ()

+ + 2 =. (6.3.13)

(r) r (r) r ()

Now we observe that keeping r fixed and varying , we find that necessarily the expression

()

()

must be constant, i.e., () is a solution of the differential equation

() = () in (0, 2) (6.3.14)

for some real constant ; on the other hand, should be 2-periodic. Hence, all admissible

have the form

m () = eim for arbitrary m Z ,

and the corresponding in (6.3.14) is given by

m = m2 .

(r) 1 (r) m2

+ 2 =,

(r) r (r) r

or, equivalently,

r 2 (r) + r (r) + (r 2 m2 )(r) = 0 . (6.3.15)

Let us denote by m (r) any solution of such equation. The boundary conditions for m are as

follows. Since w(1, ) = 0 for all , necessarily m (1) = 0. On the other hand, when m 6= 0,

the function w(r, ) = m (r)eim would not admit a limit independent of as r 0+ , unless

m (0) = 0; so we enforce this condition. When m = 0, then w(r, ) = 0 (r) and no condition is

needed on 0 at the origin, except that of being finite as r 0+ .

Equation (6.3.15) is similar to one of the classical equations of Applied Mathematics, the Bessel

equation

x2 Y (x) + x Y (x) + (x2 2 ) Y (x) = 0 , x>0, (6.3.16)

6.4. EXPANSION IN SERIES OF EIGENFUNCTIONS 103

Figure 6.2: Plots of the Bessel functions of the first kind J0 , J1 and J2

where is an arbitrary real or complex number. The solutions of (6.3.15) with the boundary

conditions described above can be related to certain solutions of the Bessel equation, precisely to

the Bessel functions of the first kind J (x), with a non-negative integer. The behavior of these

functions for = 0, 1, 2 is shown in Fig. 6.2 (taken from Wikipedia); it is apparent that near the

origin they behave as we want that 0 , 1 , 2 behave. Indeed, we have J (0) = 0 for all integer

> 0, whereas J0 (0) is finite and non-zero.

Each function J exhibits an oscillatory behavior around the horizontal axis as x increases,

with a monotonically increasing sequence of strictly positive, simple zeroes x,k , k = 1, 2, . . . .

Thus, for any m Z and any integer k 1, let us define

Then, using (6.3.16) with x = x|m|,k r and = |m|, it is easily seen that m,k is a solution of

(6.3.15) if the constant = m,k is defined as

Summarizing, for any m Z and any integer k 1, the function wm,k defined in polar

coordinates as

wm,k (r, ) = m,k (r)m () (6.3.19)

is an eigenfunction of Problem 6.1.1 in the unit circle, with corresponding eigenvalue given by

(6.3.19). One can show that these are indeed the totality of the eigenfunctions of this problem,

which therefore is completely solved.

The knowledge of the eigenvalues and eigenfunctions of Problem 6.2.1 allows us to solve by series

the variational problem associated with the bilinear form a, namely:

Problem 6.4.1. Given f H, find u V such that

104 CHAPTER 6. SPECTRAL THEORY FOR ELLIPTIC SELF-ADJOINT PROBLEMS

This means that we are able to expand the solution u in the eigenfunction series (6.2.11), where

the (generalized) Fourier coefficients of u are expressed in terms of those of the data f . The precise

result is as follows.

X

f= fk wk

k=1

be the expansion of f in the series of eigenfunctions of the bilinear form a. Then, the Fourier

coefficients of the solution u of Problem 6.4.1 are given by

1

uk = fk , k = 1, 2, . . . , (6.4.2)

k

X 1

u= fk wk (6.4.3)

k

k=1

X

X

a(u, w ) = a( uk wk , w ) = uk a(wk , w ) = (f, w ) ,

k=1 k=1

where the second equality is justified by Proposition 6.2.9. Recalling (6.2.13), we obtain

u = f ,

(

u = f in (0, L) ,

(6.4.4)

u(0) = u(L) = 0 .

Assume that right-hand side is the constant function f = 1, so that the solution is the parabola

u(x) = 12 x(L x).

6.4. EXPANSION IN SERIES OF EIGENFUNCTIONS 105

Let us first compute the Fourier coefficients of f with respect to the eigenfunctions wk intro-

duced in Sect. 6.3.1. Since

Z L 0 for k even ,

sin(k x) dx = 2L (6.4.5)

0 L for k odd ,

k

we get

0 for k even ,

fk = 2L 2 r

for k odd .

k L

The Fourier coefficients of the solution are given by (6.4.2), i.e.,

0 for k even ,

1 r

uk = fk = 2L3 2

k

3 3

for k odd ,

k L

so that the eigenfunction expansion of the solution is as follows:

X

4L2 X 1

u(x) = uk wk (x) = sin (2m + 1) x .

3 m=0 (2m + 1)3 L

k=1

(

u = f in ,

(6.4.6)

f = 0 on ,

with = (0, L)2 . Assume as above that right-hand side is the constant function f = 1; recall

that this problem has been already considered in Sect. 4.5 while discussing the regularity of the

solution of an elliptic problem in a domain with corners.

Let us first compute the Fourier coefficients of f with respect to the eigenfunctions whk intro-

duced in Sect. 6.3.2. Recalling (6.4.5), we have

Z L Z L

0 for h or k even ,

2

fhk = sin(h x) dx sin(k y) dy = 8L2 1

L 0 L 0 L

for h and k odd .

2 hk

The Fourier coefficients of the solution are given by (6.4.2), i.e.,

0 for h or k even ,

1

uhk = fhk = 8L3 1

hk

4 for h and k odd .

hk(h2 + k2 )

so that the eigenfunction expansion of the solution is as follows:

X

X

u(x) = uhk whk (x)

h=1 k=1

2 X X

16L 1

= sin (2m + 1) x sin (2n + 1) y .

4 2 2

(2m + 1)(2n + 1)((2m + 1) + (2n + 1) ) L L

m=0 n=0

106 CHAPTER 6. SPECTRAL THEORY FOR ELLIPTIC SELF-ADJOINT PROBLEMS

6.5 Exercises

6.1. Compute the eigenvalues and the eigenfunctions of the one-dimensional boundary-value prob-

lems (6.3.4) and (6.3.5).

u = xy 2 in ,

u = 0 on ,

u = x sin 3y in

u = 0 on 1 3

u

= 0 on 2 4

n

where = (0, )2 IR2 , 1 = (0, ){0}, 2 = {}(0, ), 3 = (0, ){} and 4 = {0}(0, ).

Chapter 7

Parabolic Problems

Parabolic problems describe propagation phenomena with infinite speed, the so-called diffusion

phenomena. Many problems in the applied sciences lead to this type of mathematical model: the

heat propagation through a rod or the motion of a viscous flow in a channel are just two important

examples, from Thermodynamics and Fluid Dynamics, in which parabolic equations describe the

time evolution of temperature and velocity, respectively.

Initial/boundary-value problems for a second-order linear parabolic operator can be studied in

a manner similar to what has been done in the previous chapters for elliptic problems. Precisely,

we first obtain a weak, or variational, formulation of the problem, by proceeding initially in a

formal manner and then making the functional assumptions fully rigorous. Next, we prove the

well-posedness of the weak formulation; in the present situation, the result will be a consequence

of an a-priori bound on any possible solution of the weak problem. At last, we interpret the weak

solution as a strong, or classical, one, provided suitable assumptions on the data are satisfied.

While the variational treatment of the spatial part of the operator follows the guidelines es-

tablished in Chaps. 3 and 4, here the mathematical novelty is represented by the first-order time

derivative. Its treatment will be based on a new result, that represent a generalization of the

well-known integration-by-part formula in one dimension.

Before starting, let us introduce some slightly new notations that are needed in order to take

into account the time variable.

Let be a bounded open set in IRn and suppose its boundary is smooth enough; further-

more, let (0, T ) be the time interval of interest. We introduce the cylinder in space and time

Q = (0, T ) IRn+1 .

Consider a generic real-valued function v = v(x , t) defined on Q; it is convenient to think of

it as a function v = v(t) defined in for every fixed t (0, T ) (i.e., a function in depending on

t as a parameter). In other words, for all t (0, T ), the function v(t) : IR is defined as

(v(t))(x ) = v(x , t) x .

Z Z T Z

kvk2L2 (Q) = 2

|v(x , t)| dx dt = |v(x , t)|2 dx dt

Q 0

kv(t)k2L2 () = |v(x , t)|2 dx ,

107

108 CHAPTER 7. PARABOLIC PROBLEMS

it results Z T

kvk2L2 (Q) = kv(t)k2L2 () dt .

0

Since the left-hand side is, by assumption, a finite number, it follows that we can think of v as

a function v : (0, T ) L2 () such that the further function t 7 kv(t)kL2 () belongs to L2 (0, T ).

We write v L2 (0, T ; L2 ()).

Thus, we may regard functions in L2 (Q) as functions of the variable t taking values in the

space L2 (), with the further property that the spatial norm has a certain degree of integrability

with respect to the time. This can be generalized as follows.

Lp (0, T ; X)

Z T

kv(t)kpX dt < + ,

0

Z T 1/p

kvkLp (0,T ; X) = kv(t)kpX dt

0

(with the obvious change when p = ). In addition, Lp (0, T ; X) is a Banach space if and only if

X is so.

Hence, we can identify L2 (Q) with the space L2 (0, T ; L2 ()); in this and next chapters, we

shall also use spaces like L2 (0, T ; H10 ()) and L2 (0, T ; H1 ()). We also define C0 ([0, T ]; X) as

the space of all continuous functions v : [0, T ] X, equipped with the norm

0tT

In order to illustrate how to derive the variational formulation of a parabolic problem, we put

ourselves in the conceptually simplest situation, namely we consider the homogeneous Dirichlet

problem for the heat equation:

u

u = f in Q

t (7.1.1)

u = 0 on (0, T )

u = u0 on {0}

where u = u(x , t) is the unknown temperature, at position x and time t, of a conducting body

occupying the domain , f = f (x , t) and u0 = u0 (x ) are two given functions representing the

heat exchange with the surrounding environment and the initial temperature, respectively. This

choice is motivated by the arguments of Sect. 1.4, where it has been shown that the heat operator

is the canonical form to which we can reduce any second-order parabolic operator. However, we

will also mention extensions of the subsequent results to more general situations.

7.1. VARIATIONAL FORMULATION 109

As we know, the weakest way to give meaning to the heat equation is in the distributional

sense: we assume u and f to be locally integrable in Q and we require u to satisfy

u

u = f in D (Q) , (7.1.2)

t

i.e., Z Z Z

u dx dt u dx dt = f dx dt D(Q) . (7.1.3)

Q t Q Q

A more balanced formulation, suggested by the experience gained with elliptic problems, consists

of applying Gauss therem only once in the term involving the spatial operator, i.e., writing the

term Z Z

u dx dt as u dx dt .

Q Q

Obviously, this requires more regularity on u: it is quite natural to assume u L2 (0, T ; H10 ()),

since then also the Dirichlet boundary condition is rigorously defined a.e. (almost everywhere, i.e.,

outside a set of zero measure) in time. Furthermore, this assumption allows us to give a precise

meaning to the time derivative of u. Indeed, recalling the bound (3.9.11), one has

Z T Z T

kuk2H1 () dt C kuk2H1 () dt < + ,

0

0 0

i.e.,

u L2 (0, T ; H10 ()) u L2 (0, T ; H1 ()) .

If we furtherly assume that f L2 (Q) = L2 (0, T ; L2 ()) L2 (0, T ; H1 ()) (the latter inclusion

being a consequence of (3.9.1)), then we deduce from (7.1.2) that

u

= u + f L2 (0, T ; H1 ()) . (7.1.4)

t

Thus, under the above assumption on f , we are led to look for a solution u of problem (7.1.1)

in the space

w

W(0, T ; H10 (), H1 ()) = {w L2 (0, T ; H10 ()) : L2 (0, T ; H1 ())} , (7.1.5)

t

where the time derivative has to be understood, as usual, in the distributional sense. This is a

Hilbert space for the graph norm

1/2

w 2

kwkW(0,T ; H10 (),H1 ()) = kwk2L2 (0,T ; H1 ()) +k k 2 1 .

0 t L (0,T ; H ())

In this case, all three addends of the heat equation are in L2 (0, T ; H1 ()), and the following

variational formulation of the equation can be given:

Z T Z T

u

1

H () h u, vi 1

H0 () dt = H1 () hf, viH10 () dt v L2 (0, T ; H10 ()) . (7.1.6)

0 t 0

What about the initial condition? We have to give a correct meaning to it. This is precisely

what is provided by the following result, which in addition establishes a useful integration-by-

parts formula in time for functions in W(0, T ; H10 (), H1 ()).

110 CHAPTER 7. PARABOLIC PROBLEMS

Proposition 7.1.1. Any function in W(0, T ; H10 (), H1 ()) is (up to a negligible set) con-

tinuous from [0, T ] to L2 (); in other words, the space W(0, T ; H10 (), H1 ()) is contained in

C0 ([0, T ]; L2 ()) with continuous inclusion.

Furthermore, for any w, v W(0, T ; H10 (), H1 ()), the following identity holds:

d w v

(w, v)L2 () = H1 () h , viH10 () + H1 () h , wiH10 () in D (0, T ) ; (7.1.7)

dt t t

equivalently, for any 0 t1 < t2 T , one has

Z t2 Z t2

w v

H1 () h , viH10 () dt + H1 () h , wiH10 () dt = (w(t2 ), v(t2 ))L2 () (w(t1 ), v(t1 ))L2 () .

t1 t t1 t

(7.1.8)

Proof. (We just provide the essential steps.) At first, one proves that the set of all smooth enough

functions, e.g. those in C1 (Q) vanishing on [0, T ], are dense in W(0, T ; H10 (), H1 ()); we

skip the technical details.

Next, we prove (7.1.7). Let {vn }n0 be any sequence of smooth functions converging to v in

W(0, T ; H10 (), H1 ()). Then, for all D(0, T ), one has

Z T Z T

d

D (0,T ) h (w, v)L2 () , iD(0,T ) = (w(t), v(t))L 2 () (t) dt = lim (w(t), vn (t))L2 () (t) dt .

dt 0 n 0

Now,

Z T Z T Z

d

(w(t), vn (t))L2 () (t) dt = w(x , t)v(x , t)

(t) dx dt

0 0 dt

Z T Z Z TZ

vn

= w(x , t) vn (x , t)(t) dx dt w(x , t)(t) (x , t) dx dt .

0 t 0 t

w

By definition of t L2 (0, T ; H1 ()), observing that vn D(Q), we can write

Z T Z Z T

w

w(x , t) vn (x , t)(t) dx dt = H1 () h , vn iH10 () dt .

0 t 0 t

Now, we remember that - since (H10 (), L2 (), H1 ()) form a Gelfand triple (see Sect. 3.9), the

L2 ()-inner product of a function g L2 () with a function z H10 () can be equivalently viewed

as the duality pairing between H1 () and H10 (), i.e.,

Therefore,

Z T Z Z T

vn vn

w(x , t)(t) (x , t) dx dt = H1 () h , wiH10 () dt .

0 t 0 t

Summarizing,

Z T Z T Z T

w vn

(w(t), vn (t))L2 () (t) dt = H1 () h , vn iH10 () dt + H1 () h , wiH10 () dt .

0 0 t 0 t

7.1. VARIATIONAL FORMULATION 111

Letting n , we get

Z T Z T

d w v

D (0,T ) h (w, v)L2 () , iD(0,T ) = H1 () h , viH10 () dt + H1 () h , wiH10 () dt

dt 0 t 0 t

Z T

w v

= 1

H () h , vi 1

H0 () + 1

H () h , wiH0 () dt

1

0 t t

w v

= D (0,T ) h H1 () h , viH10 () + H1 () h , wiH10 () , iD(0,T ) .

t t

Thus, (7.1.7) is proven.

This relation indicates that the L1 (0, T )-function t 7 (w(t), v(t))L 2 () has first derivative, in

the distributional sense, which itself belongs to L1 (0, T ). From such a result one can easily deduce

(as in the proof of Property 3.2.7) that

the function (w(t), v(t))L 2 () is continuous (indeed, absolutely continuous) on [0, T ] . (7.1.10)

In order to prove the continuity in L2 () of w(t) at any point t0 [0, T ], we use (7.1.10) once

with v = w and then with the constant function v = w(t0 ). Using the identity

kw(t) w(t0 )k2L2 () = (w(t), w(t)L 2 () 2(w(t), w(t0 ))L2 () + (w(t0 ), w(t0 )2L2 ()

At last, (7.1.8) is obtained by integrating (7.1.7) on the time interval [t1 , t2 ].

Proposition 7.1.1 suggests the choice of the initial data u0 in L2 (), since u(0) is well-defined

as a function in this space, if u W(0, T ; H10 (), H1 ()); thus, condition u(0) = u0 has to be

understood as an equality between functions in L2 ().

Remark 7.1.2. Proposition 7.1.1 is just a particular case of a more general result, due to J.-

L. Lions, which concerns Gelfand triples (V, H, V ) (where V H with continuous and dense

inclusion). Precisely, introducing the space

w

W(0, T ; V, V ) = {w L2 (0, T ; V ) : L2 (0, T ; V )} , (7.1.11)

t

one has:

Proposition 7.1.3. The space W(0, T ; V, V ) is contained in C0 ([0, T ]; H) with continuous in-

clusion. Furthermore, for any w, v W(0, T ; V, V ) and any 0 t1 < t2 T , one has

Z t2 Z t2

w v

V h

, viV dt + V h , wiV dt = (w(t2 ), v(t2 ))H (w(t1 ), v(t1 ))H . (7.1.12)

t1 t t1 t

Note that if one takes H = V (hence, V = H = V ) the result tells us that if w L2 (0, T ; V )

w

and L2 (0, T ; V ), then w C0 ([0, T ]; V ).

t

We are ready for stating a variational formulation of Problem (7.1.1). For simplicity, in the

sequel we will set

Z Z

a(u, v) = (u, v)H10 () = u v dx , (u, v) = (u, v)L2 () = uv dx .

112 CHAPTER 7. PARABOLIC PROBLEMS

Problem 7.1.4. Given f L2 (Q) and u0 L2 (), find u W(0, T ; H10 (), H1 ()) satisfying

u(0) = u0 and such that

Z T Z T Z T

u

H1 () h , viH10 () dt + a(u, v) dt = (f, v) dt v L2 (0, T ; H10 ()) . (7.1.13)

0 t 0 0

Two alternative but entirely equivalent expressions of (7.1.13) can be given. They are:

u

H1 () h (t), viH10 () + a(u(t), v) = (f (t), v) v H10 (), a.e. in (0, T ) (7.1.14)

t

and

d

(u(t), v) + a(u(t), v) = (f (t), v) v H10 (), a.e. in (0, T ). (7.1.15)

dt

Expression (7.1.15) is particularly important both for proving the existence of a solution (as we

shall see later) and for the numerical approximation of the problem. The forms (7.1.13) and (7.1.14)

are equivalent, thanks to the fact that the set of all piecewise constant functions v : [0, T ] H10 ()

is dense in L2 (0, T ; H10 ()). On the other hand, (7.1.14) and (7.1.15) are equivalent, since

d u

(u(t), v) = H1 () h (t), viH10 () in D (0, T ) ,

dt t

thanks to (7.1.7) with v L2 (0, T ; H10 ()) constant in time.

In this section, we establish an upper bound for certain norms of any solution u of the variational

problem (7.1.4), depending only on suitable norms of the data f and u0 . Such an estimate is

called a priori because we derive it from the sole assumption of u being a solution of this problem,

without even knowing that such a function exists.

However, a result of this kind is of great importance since from it we may prove the well-

posedness of our parabolic problem, i.e., the existence and uniqueness of the solution as well as

its continuous dependence on the data.

Let us suppose that u is a solution of Problem (7.1.4). Fix any time satisfying 0 < T .

Take v = u in (7.1.14) and integrate in time from 0 to (equivalently, in (7.1.13) choose v

coinciding with u in [0, ] and equal to 0 in (, T ]); we get

Z Z Z

u

1

H () h , ui 1

H0 () dt + a(u, u) dt = (f, u) dt . (7.2.1)

0 t 0 0

The first integral can be expressed via Proposition 7.1.1 as follows. Set t1 = 0, t2 = and

w = v = u in (7.1.8); this easily yields

Z

u 1

H1 () h , uiH10 () dt = [(u( ), u( )) (u(0), u(0))]

0 t 2

1 1

= ku( )k2L2 () ku0 k2L2 () .

2 2

On the other hand, Z Z

a(u, u) dt = kuk2H1 () dt .

0

0 0

7.2. AN A PRIORI ESTIMATE 113

1

where CP () is the Poincare constant. Finally, using the inequality ab 2 a2 + b2 a, b 0

with the choice a = CP () kf kL2 () and b = kukH10 () leads to

1 2 1

(f, u) CP () kf k2L2 () + kuk2H1 () ;

2 2 0

Z Z Z

1 1

(f, u) dt CP2 () kf k2L2 () dt + kuk2H1 () dt .

0 2 0 2 0 0

Z Z Z

1 2 2 1 2 1 2 2 1

ku( )kL2 () + kukH1 () dt ku0 kL2 () + CP () kf kL2 () dt + kuk2H1 () dt ,

2 0 0 2 2 0 2 0 0

Z Z

2 2 2 2

ku( )kL2 () + kukH1 () ku0 kL2 () + CP () kf k2L2 () dt . (7.2.2)

0

0 0

Now, if we neglect the second term on the left-hand side and use the trivial bound

Z Z T

2

kf kL2 () dt kf k2L2 () dt ,

0 0

we get Z T

ku( )k2L2 () ku0 k2L2 () + CP2 () kf k2L2 () dt for all [0, T ] ,

0

i.e.,

Z T

max ku( )k2L2 () ku0 k2L2 () + CP2 () kf k2L2 () dt .

[0,T ] 0

On the other hand, if we neglect the first term on the left-hand side of (7.2.2) and we choose

= T , we have Z T Z T

2 2 2

kukH1 () ku0 kL2 () + CP () kf k2L2 () dt .

0

0 0

p

Taking the square roots of both sides in the two last inequalities, and using the relation 2 + 2

+ for all , 0, we end up with the following result.

Proposition 7.2.1. Any solution u of Problem 7.1.4 satisfies the estimate

kukC0 ([0,T ]; L2 ()) + kukL2 (0,T ; H10 ()) C ku0 kL2 () + kf kL2 (0,T ; L2 ()) , (7.2.3)

Remark 7.2.2. A similar result holds, with the obvious changes, if f L2 (0, T ; H1 ()), or if

a(u, v) is any bilinear form which is coercive in H10 () (as the one associated with a uniformly

elliptic second-order operator in ).

114 CHAPTER 7. PARABOLIC PROBLEMS

We now derive several important consequences from the a priori estimate (7.2.3).

(i) Suppose that two sets of data {f1 , u0,1 } and {f2 , u0,2 } are given and denote by u1 and u2

the solutions of the corresponding Problems 7.1.4. Due to the linearity of the equations,

u1 u2 is the solution of the problem whose data are f1 f2 and u0,1 u0,2 , so that from

(7.2.3) we have

(7.3.1)

C ku0,1 u0,2 kL2 () + kf1 f2 kL2 (0,T ; L2 ()) ,

which shows that as the data set {f1 , u0,1 } approaches {f2 , u0,2 }, the first solution u1 also

approaches the second solution u2 .

This means that the solution of the problem, if it exists, depends continuously on the data.

(ii) Consider a set of data {f, u0 } and suppose that two solutions u1 and u2 arise; then we must

have

ku1 u2 kC0 ([0,T ]; L2 ()) + ku1 u2 kL2 (0,T ; H10 ()) 0

and thus u1 = u2 . This means that the solution of the problem, if it exists, is unique.

Let {n , wn }n1 be the eigenvalue-eigenfunction pairs of the Laplacian with Dirichlet boundary

conditions, given by the problem (??); since {wn }n1 is an orthonormal fundamental set in L2 (),

we can give the following representations for the functions f (t), u0 L2 ():

+

X

f (t) = fn (t)wn , fn (t) = (f (t), wn )

n=1

+

X

u0 = u0,n wn , u0,n = (u0 , wn ).

n=1

+

X

u(t) = un (t)wn ; (7.3.2)

n=1

substituting into equation (7.1.15) and choosing any eigenfunction wm as a test function, gives

+

! +

! +

!

d X X X

un (t)wn , wm + a un (t)wn , wm = fn (t)wn , wm , m = 1, 2, . . . ,

dt

n=1 n=1 n=1

that is

+

X +

X +

X

un (t) (wn , wm ) + un (t) a(wn , wm ) = fn (t) (wn , wm ) m = 1, 2, . . . ;

n=1 n=1 n=1

recalling that

(wn , wm ) = n,m and a(wn , wm ) = n n,m ,

7.3. WELL-POSEDNESS OF THE PROBLEM 115

we obtain

um (t) + m um (t) = fm (t) , m = 1, 2, . . . (7.3.3)

which shows that every generalized Fourier coefficient um (t) of the expansion of u satisfies a

first-order linear ordinary differential equation.

Moreover, since u(0) = u0 , the following relationship must hold:

+

X +

X

un (0)wn = u0,n wn ,

n=1 n=1

thus

um (0) = u0,m m = 1, 2, . . . . (7.3.4)

Finally, from (7.3.3) and (7.3.4) we have for every m 1 the Cauchy problem

(

um (t) + m um (t) = fm (t) ,

(7.3.5)

um (0) = u0,m ,

whose solution is Z t

um (t) = e m t

u0,m + em ( t) fm ( ) d. (7.3.6)

0

It remains now to be checked that the series (7.3.2) converges to a function u which solves

Problem 7.1.4. To do this, let us set

N

X N

X N

X

uN (t) : = un (t)wn , fN (t) : = fn (t)wn , u0,N : = u0,n wn ;

n=1 n=1 n=1

it is easy to verify that uN satisfies the variational formulation of the problem

uN

uN = fN in Q ,

t

uN = 0 on (0, T ) ,

uN = u0,N on {0} ;

then, if we consider a further set of data {fM , u0,M } with the corresponding solution uM we have,

from the upper bound (7.2.3),

C ku0,N u0,M kL2 () + kfN fM kL2 (0,T ; L2 ()) , M, N 1 . (7.3.7)

Since {fN }N 1 and {u0,N }N 1 converge in their spaces (to f and u0 , respectively), they are

Cauchy sequences; from (7.3.7) it follows that so is {uN }N 1 both in C0 ([0, T ]; L2 ()) and in

L2 (0, T ; H10 ()). Using the completeness of these spaces allows us to conclude that when N

the sequence uN converges to a function u belonging to L2 (0, T ; H10 ()) C0 ([0, T ]; L2 ()).

Passing to the limit in the distributional equations satisfied by uN , i.e.,

Z Z Z

uN dx dt uN dx dt = fN dx dt D(Q) ,

Q t Q Q

we obtain (7.1.2), from which one can easily deduce that u solves Problem 7.1.4.

Summarizing, we have proven the following fundamental result.

Theorem 7.3.1. Problem 7.1.4 admits one and only one solution, for which the bound (7.2.3)

holds. Furthermore, the solution depends continuously on the data, as expressed by the bound

(7.3.1).

116 CHAPTER 7. PARABOLIC PROBLEMS

We have seen that, under the hypothesis f L2 (0, T ; L2 ()) and u0 L2 (), the parabolic

problem (7.1.1) admits a unique solution u L2 (0, T ; H10 ()); note that we need not assign u0 in

H10 (): u0 L2 () is enough to obtain a solution which belongs to H10 () a. e. in time. This fact

is already related to the regularization property of parabolic equations.

Another point of view is provided by a spectral analysis. Let us suppose f 0, i.e., fn (t) 0

for every n 1; then, the solution reads as

+

X +

X

u(t) = un (t)wn = en t u0,n wn .

n=1 n=1

P+

Recalling that the L2 -norm

P+ of a2 function v of the form v = n=1 vn wn is given by Parsevals

2

identity as kvkL2 () = n=1 |vn | , it results

+

X

ku(t)k2L2 () = e2n t |u0,n |2 ;

n=1

+

X

|u0,n |2 = ku0 k2L2 () < +

n=1

|un (t)| Cen t n 1 .

It follows that the Fourier coefficients of the solution for any time t > 0 decay exponentially in

the wavenumber n. As noted in Chap. 6, this is a manifestation of smoothness of the solution.

If f is not zero, an analogous conclusion can be achieved from expression (7.3.6). Since t < 0

whenever t > 0, any singularity of the data f at a given time > 0 is immediately smoothed out

at all subsequent times.

In short, a parabolic problem tends to smooth the data.

Another explicit manifestation of this property can be obtained by considering the homoge-

neous heat equation in the half-plane {(x, t) : t > 0}:

u 2 u

2 = 0 in IR (0, +) ,

t x

(7.4.1)

u(x, t) 0 for |x| ,

u(x, 0) = u0 (x) x IR .

If we apply the Fourier transform with respect to x to u and its derivatives we get

Z +

1

u(, t) = u(x, t)eix dx ,

2

d

2u

(, t) = 2 u(, t) ,

x2

c

u u

(, t) = (, t) ,

t t

7.4. SOME FACTS ABOUT THE REGULARITY OF THE SOLUTION 117

and then

u + 2 u = 0

t IR, t > 0 . (7.4.2)

u(, 0) = u0 ()

Note that the condition u(x, t) 0 is implicit in the fact that u(t) admits a Fourier transform,

i.e., u(t) L2 (IR) and so it need not be furtherly imposed in Problem (7.4.2).

Solving this Cauchy problem with respect to t and regarding as a parameter yields

2

u(, t) = u0 ()e t ;

Z + Z +

1 ix 1 2

u(x, t) = u(, t)e d = u0 ()e t+ix d =

2 2

Z + Z +

1 2

= e t+ix u0 (y)eiy dy d =

2

Z + Z +

1 2 t+i(xy)

= e d u0 (y) dy.

2

If we set Z +

1 2 t+i(xy)

K(x y, t) : = e d , (7.4.3)

2

we can write the solution as

Z +

u(x, t) = K(x y, t)u0 (y) dy = (K u0 )(x, t) ,

where the convolution is intended with respect to the space variable only.

The function K is said to be the kernel of the heat equation; it is the solution of the parabolic

problem

K 2K

= 0 in IR (0, +) ,

t x2

K 0 for |x| ,

K = 0 for t = 0 ,

in the sense of distributions (for this reason it is sometimes called a fundamental solution of the

heat equation) and it can be expressed in a closed form, since from equation (7.4.3) it follows

z2

e 4t

K(z, t) = ,

4t

so that we also have

Z +

(xy)2

e 4t

u(x, t) =

u0 (y) dy .

4t

From here we see that u is infinitely differentiable with respect to x even when u0 is not, since

the kernel K lies in C for every t > 0.

An alternative fashion for studying the regularity of the solution of a parabolic problem consists

of the following procedure:

118 CHAPTER 7. PARABOLIC PROBLEMS

u

u = f

t

and formally differentiate both left-hand side and right-hand side with respect to t to obtain

an equation satisfied by u

t ;

u

u = f t > 0

t

and exploit the results about the regularity of the solution of elliptic problems.

u u f

= (7.4.4)

t t t t

f

so u

t formally satisfies a parabolic equation where the right-hand side is t . Furthermore, from

u = 0 on (0, T ), it follows the boundary condition

u

= 0 on (0, T ) ,

t

and from u(0) = u0 on {0} we obtain analogously

u

(0) = f (0) + u(0) = f (0) + u0

t

giving us the initial condition for equation (7.4.4).

In order to justify the previous formal steps, let us suppose f 2 2

t L (0, T ; L ()): since we

already assumed f L2 (0, T ; L2 ()), we derive f C0 ([0, T ]; L2 ()) from Remark 7.1.2, so that

f (0) is well-defined in L2 (); moreover, let u0 H2 (), and consequently u0 L2 (). Therefore,

the problem

w f

t w = t in Q ,

w = 0 on (0, T ) ,

w = f (0) + u0 on {0} ,

admits a unique solution w L2 (0, T ; H10 ()) with w 2 1

t L (0, T ; H ()). But it is clear that

u

one has w = t (at least in the sense of distributions) and then u 2 1

t L (0, T ; H0 ()) with

2u

t2

L2 (0, T ; H1 ()), i.e. by Proposition 7.1.1, u 0 2

t C ([0, T ]; L ()).

As a second step, we write

u

u(t) = f (t) (t) a. e. in t , (7.4.5)

t

which is an elliptic equation with respect to x , whose right-hand side lies in L2 () for almost

every t (0, T ). If is smooth or convex, then the following relationship holds true:

u

ku(t)kH2 () C
f (t) (t)
< + , C>0,

t
L2 ()

We conclude that the solution is more regular than in the general case, thanks to the extra

regularity assumptions made on the data f and u0 .

7.5. THE MAXIMUM PRINCIPLE FOR PARABOLIC EQUATIONS 119

We can extend the Maximum Principle presented in Chapter 5 for elliptic equations to parabolic

equations of the form

u

+ Lu = f in Q

t

u = g on (0, T )

u = u0 on {0}.

where L denotes the most general second order linear elliptic operator (cf. equation (??)).

By the same variational technique of Stampacchia Theorem 5.2.2, and using now the a priori

upper bound of the parabolic problems, it is possible to prove the following result:

m = min min g, min u0 and M = max max g, max u0

and assume that the same hypothesis of Theorem 5.2.2 hold on L. Then:

(i) if f a0 m 0 in Q, then u m in Q;

(ii) if f a0 M 0 in Q, then u M in Q.

7.6 Exercises

7.1. Solve in series of eigenfunctions the following parabolic problem:

u 2 u

2 = 0 in (0, 1) (0, T ), T > 0 ,

t x

u(0, t) = et o

0<tT ,

u(1, t) = 0

u(x, 0) = 0 x (0, 1) .

(Hint: do a suitable substitution to have a homogeneous Dirichlet condition on the boundary; for

example, consider u(x, t) = u(x, t) + et (1 x) and solve first for u).

u 2 u

2 = x in (0, ) (0, T ), T > 0 ,

t

x

)

u(0, t) = 0

u 0<tT ,

(, t) = 0

n

u(x, 0) = 1 x (0, ) .

(Hint: first of all, note that in this case it results n x . Moreover, for this problem you need

the eigenfunctions of the Laplacian with mixed boundary conditions: see Exercise (6.1)).

120 CHAPTER 7. PARABOLIC PROBLEMS

Chapter 8

Hyperbolic Problems

Hyperbolic problems arise in modeling transport phenomena with finite speed, especially those

involving wave motion. Maxwells equations, which describe the propagation of electromagnetic

waves in the vacuum, as well as DAlemberts equation for the pressure waves in a fluid like air or

water, and the so-called equation of a vibrating string, which gives the motion of an elastic wave

induced across a rope blocked at its ends, go all back to a hyperbolic mathematical model.

In this chapter we deal with the most general second order linear hyperbolic operator in its

canonical form, as it has been presented in Chapter 1. We shall use the same notation which has

been introduced in the previous chapter and our explanation will be almost entirely parallel to

that of parabolic problems.

Throughout this chapter we shall refer to the initial/boundary value problem

2

u

u = f in Q ,

2

t

u = 0 ) on (0, T ) , (8.1.1)

u = u0

u on {0} .

= u1

t

Note that in this case we must prescribe two initial conditions due to the presence of the second

2

order derivative t2u in time.

2u

c2 u = f , (8.1.2)

t2

where c is the speed of the wave which is propagating across the domain .

However, it is possible to reduce such an equation to the form discussed in the text by a classical

procedure of the Mathematical Physics which consists in rewriting the problem in a dimensionless

form.

For instance, suppose u is a velocity field and let us denote by U a characteristic (constant)

velocity such that u = U u, where u is a dimensionless variable; analogously, let us set x = Lx and

121

122 CHAPTER 8. HYPERBOLIC PROBLEMS

t = t, L and being a characteristic length of the spatial domain and a characteristic time,

respectively. Substituting into equation (8.1.2) gives

U 2 u c2 U

2 u = f ,

2 t2 L

that is

2 u c2 2 2

u = f.

t2 L2 U

Choosing = L/c and setting 2 f /U = f, we finally obtain

2 u

u = f ,

t2

which is a dimensionless equation of the desired form, in the unknown u.

2u

u = f (8.1.3)

t2

certainly makes sense in D (Q), so we can write

2u

D (Q) h u, iD(Q) = D (Q) hf, iD(Q) D(Q)

t2

or, more explicitly,

Z T Z T Z T

2

u, dt + a(u, ) dt = (f, ) dt D(Q), (8.1.4)

0 t2 0 0

Z Z

a(u, v) = (u, v)H10 () = u v dx , (u, v) = (u, v)L2 () = uv dx .

On the other hand, if we assume f L2 (Q) = L2 (0, T ; L2 ()) and u L2 (0, T ; H10 ()), which

implies u L2 (0, T ; H1 ()), we obtain from (8.1.3)

2u

L2 (0, T ; H1 ()) ,

t2

so that we can express our equation in the following variational form:

Z T Z T Z T

2u

H1 () h , viH10 () dt + a(u, v) dt = (f, v) dt v L2 (0, T ; H10 ()) . (8.1.5)

0 t2 0 0

In order to give a precise meaning to the initial conditions, we make the further assumption

that

u

L2 (0, T ; H10 ()) ;

t

in other words, all together the solution u is required to satisfy u L2 (0, T ; H10 ()) with u

t

1 1 u 2 u

W(0, T ; H0 (), H ()). As a consequence, Proposition 7.1.1 applied to the pair t , t2 yields

u

C0 ([0, T ]; L2 ()) ;

t

8.2. AN A PRIORI ESTIMATE 123

u

on the other hand, recalling Remark 7.1.2 for the pair u, t , we get

u

Thus, the pointwise evaluation of both u and t makes sense in [0, T ]. In particular, we have

u

u(0) = u0 H10 () and (0) = u1 L2 (),

t

giving the natural spaces in which one has to choose the initial data.

We are ready to state the initial/boundary value problem (8.1.1) in variational form.

Problem 8.1.2. Given f L2 (Q), u0 H10 () and u1 L2 (), find u L2 (0, T ; H10 ()) with

u 1 1 u

t W(0, T ; H0 (), H ()) satisfying u(0) = u0 , t (0) = u1 and such that

Z T Z T Z T

2u

1

H () h , viH10 () dt + a(u, v) dt = (f, v) dt v L2 (0, T ; H10 ()) . (8.1.6)

0 t2 0 0

2u

H 1

() h (t), viH10 () + a(u(t), v) = (f (t), v) v H10 (), a.e. in (0, T ) (8.1.7)

t2

and

d2

(u(t), v) + a(u(t), v) = (f (t), v) v H10 (), a.e. in (0, T ). (8.1.8)

dt2

As we have already seen for parabolic problems, obtaining an a priori upper bound on the norm of

the solution u is once again of very great importance to prove the well-posedness of the problem

(8.1.1).

Consider any (0, T ]. Taking v = u

t (t) in (8.1.7) and integrating from 0 to yields

Z Z Z

2 u u u u

H 1

() h , i 1 dt + a(u, ) dt = (f, ) dt . (8.2.1)

0 t2 t H0 () 0 t 0 t

2u u

Now take t1 = 0, t2 = , w = t2 and v = t in (7.1.7) and use the initial condition to get

Z

2 u u 1 u u u u

H1 () h 2

, iH10 () dt = ( ), ( ) (0), (0) =

0 t t 2 t t t t

2

1
u
1

=
( )
ku1 k2L2 () .

2 t 2 2

L ()

Furthermore, the following identity holds for our bilinear symmetric form a(u, v):

u 1 d

a u, = a(u, u) .

t 2 dt

124 CHAPTER 8. HYPERBOLIC PROBLEMS

1

a(u(t + t), u(t + t)) a(u(t), u(t))

t

1

= a(u(t + t), u(t + t)) a(u(t), u(t + t)) a(u(t), u(t))

t

u(t + t) u(t) u(t + t) u(t)

=a , u(t + t) + a u(t), ,

t t

Z

u 1 1 1

a u, dt = [a(u( ), u( )) a(u(0), u(0))] = ku( )k2H1 () ku0 k2H1 () .

0 t 2 2 0 2 0

Finally, using the Cauchy-Schwartz inequality and the fact ab 12 a2 + b2 a, b > 0, we get

Z Z

u 1
u

f, dt kf (t)kL2 ()
(t)

dt

0 t 0 t
L2 ()

Z Z

1
u
2

2

kf (t)kL2 () dt +
(t)
dt

2 0 2 0
t
L2 ()

These results lead us to the following relationship:

2

1

u ( )
1 1 1

ku1 k2L2 () + ku( )k2H1 () ku0 k2H1 ()

2 t L2 () 2 2 0 2 0

Z

1
u
2

kf k2L2 (0, ; L2 ()) +
(t)
dt ;

2 2
t
2

0 L ()

u
2

max
(t)
max ku(t)k2H1 ()

t[0,T ]
t
2 + t[0,T ] 0

L ()

Z T

1
u
2

kf k2L2 (0,T ; L2 ()) + ku0 k2H1 () + ku1 k2L2 () + max
(t)
dt .

0 t[0,T ]
t
L2 () 0

u

kukC0 ([0,T ]; H10 ()) +

t 0

C ([0,T ]; L2 ())

h i

C ku0 kH10 () + ku1 kL2 () + T kf kL2 (0,T ; L2 ()) , (8.2.2)

8.3. WELL-POSEDNESS OF THE PROBLEM 125

Remark 8.2.2. Observe that the right-hand side depends on T , which implies that the upper

bound becomes less and less tight as T grows to infinity. Although this is not a so serious problem

from a theoretical viewpoint, it may be unfavourable in a numerical approach, where one would

like to control and confine as much as possible the error on the numerical solution.

If we want to improve in this sense the a priori upper bound (8.2.2) we may, for instance,

control the norm of the solution u by using a different norm of the forcing term f . In particular,

since f L2 (0, T ; L2 ()) implies1 f L1 (0, T ; L2 ()), we have

Z Z

u
u

f, dt kf (t)kL2 ()
(t)
dt

t
t
2

0 0 L ()

Z

u

max
(t)
2 kf (t)kL2 () dt

t[0, ] t L () 0

u

t
0 kf kL1 (0,T ; L2 ())

C ([0, ]; L2 ())

2

1
u
1

+ kf k2L1 (0,T ; L2 ()) ,

2 t C0 ([0, ]; L2 ()) 2

Using the a priori upper bound (8.2.2) and following the guidelines of the parabolic case, one gets

the following result.

Theorem 8.3.1. Problem 8.1.2 admits one and only one solution, for which the bound (8.2.2)

holds. Furthermore, the solution depends continuously on the data in the norms which appear in

this bound.

Proof. (Sketch). Let us just detail the construction of the approximants of the solution by an

eigenfunction expansion; the rest of the arguments are similar to the parabolic case.

Let us denote, as usually, by {n , wn } the eigenvalue-eigenfunction pairs of the Laplacian with

Dirichlet boundary conditions; then we can represent the data as

+

X

f (t) = fn (t)wn , fn (t) = (f (t), wn ) ,

n=1

+

X

u0 = u0,n wn , u0,n = (u0 , wn ) ,

n=1

+

X

u1 = u1,n wn , u1,n = (u1 , wn ) ,

n=1

1

Given a set A, we recall that the following relationship holds true: L2 (A) L1 (A), provided A has a finite

Lebesgue measure; in this case, there exists a constant C > 0 such that kvkL1 (A) CkvkL2 (A) v L2 (A), which

shows that L1 -norm is weaker than L2 -norm. Nevertheless, these two norms are not equivalent, since the converse

need not be true; consider, e.g., the function v(x) = 1x on A = [0, 1].

126 CHAPTER 8. HYPERBOLIC PROBLEMS

+

X

u(t) = un (t)wn , (8.3.1)

n=1

+

X

u u

(t) = un (t)wn , un (t) = (t), wn .

t t

n=1

+

! +

! +

!

d2 X X X

un (t)wn , wm + a un (t), wm = fn (t)wn , wm , m = 1, 2, . . . .

dt2

n=1 n=1 n=1

Recalling that

(wn , wm ) = n,m and a(wn , wm ) = n n,m

we obtain

um (t) + m um (t) = fm (t), m = 1, 2, . . . ,

together with the conditions

+

X +

X

u(0) = un (0)wn = u0 = u0,n wn um (0) = u0,m m 1 ,

n=1 n=1

+

X +

X

u

(0) = un (0)wn = u1 = u1,n wn um (0) = u1,m m 1 ;

t n=1 n=1

thus, each generalized Fourier coefficient of the solution u is the solution of the following initial-

value problem for a linear second-order differential equation

um (t) + m um (t) = fm (t) ,

um (0) = u0,m , m = 1, 2, . . . ,

um (0) = u1,m

which can be rewritten in a canonical form as a first-order system by defining vm (t) = um (t)

um (t) = vm (t) ,

vm (t) = m um (t) + fm (t) ,

m = 1, 2, . . . ,

u (0) = u0,m ,

m

vm (0) = u1,m

and then introducing the vector vm (t) = (um (t), vm (t))T :

0 1 0

v (t) = vm (t) +

m m 0 fm (t)

m = 1, 2, . . .

u0,m

vm (0) =

u1,m

The solution of this differential system writes formally as

8.4. QUALITATIVE PROPERTIES OF THE SOLUTION 127

Z t

tAm u0,m ( t)Am 0

vm (t) = e + e d , (8.3.2)

u1,m 0 fm ( )

0 1

where we have set Am = . This matrix has the following eigenvalues: i m , thus it

m 0

can be diagonalized using a nonsingular matrix P CI 2,2 such that

i m 0

Am P = P ,

0 i m

which also yields the explicit expression of the exponential matrix

!

ei m t 0

etAm = P P 1 .

0 ei m t

We briefly comment on certain properties of the solution of a second-order hyperbolic problem, as

far as regularity, conservation and time reversibility are concerned.

A hyperbolic operator, unlike a parabolic one, does not smooth the solution as time increases;

roughly speaking, we may say that it simply mix together the data. This property, related to

the fact that the operator describes a propagation phenomenon, was already observed in Chap.

1, where we provided the analytical expression of the solution in the one dimensional case. The

spectral decomposition of the solution, given above, yields another point of view.

Indeed, assume that f 0, i.e., fm (t) = 0 for all m 1. Then, the general expression (8.3.2)

gives

um (t) = C1 eim t + C2 eim t ,

with m = m . The coefficients C1 , C2 depend on u0,m and u1,m but not on t. As a consequence,

if we choose, respectively, u0 and u1 to lie in L2 () and in H10 () the solution u does not become

more and more regular as the time goes by, since the terms eim t do not decay in time. In other

words, the series

+

X

sn |un (t)|2

n=1

converges in general only for s = 0 and s = 1, according to the fact that u L2 (0, T ; H10 ()).

If f is not zero, an analogous result can be achieved.

In order to have a smooth solution (in the Sobolev or classical scale) one has to assume

sufficiently smooth initial data and forcing term, plus of course the smoothness of the domain if

one is interested to the regularity up to the boundary.

Conservation properties

Let us consider an elastic membrane blocked at its boundary, which can oscillate from an initial

position u0 and with an initial speed u1 . The mathematical dimensionless model describing such

a phenomenon is

128 CHAPTER 8. HYPERBOLIC PROBLEMS

2

u

u = 0 in (0, ) ,

2

t

u = 0 ) on (0, ) , (8.4.1)

u = u 0

u on {0} ,

= u1 .

t

that is a linear hyperbolic problem with in particular the forcing term f equal to zero. Due to

this, the a priori upper bound derived from (8.2.1) can now be written as an equality:

2

1

u ( )
1 2 1 2 1 2

2 t
2 + 2 ku( )kH10 () = 2 ku1 kL2 () + 2 ku0 kH10 () , 0 .

L ()

Thus, the expression
2

1
u
1

E(t) =
(t)

+ ku(t)k2H1 ()

2 t L2 () 2 0

denotes a quantity which does not vary during the evolution in time, i.e.,

d

E(t) 0 ;

dt

In other words, the quantity E(t) is conserved in time. Physically, it represents the membranes

dimensionless total energy, and, more precisely,

2

1
u

(T )

is the dimensionless kinetic energy ,

2 t 0,

1

|u(T )|21, is the dimensionless elastic energy .

2

Instead, the following problem

2

u u

+ u = 0 in (0, ) ,

2

t

t

u = 0 ) on (0, ) ,

u = u0

u on {0} ,

= u1

t

with > 0 models the motion of the same membrane but without neglecting friction forces,

proportional to its velocity, which are taken into account by the term u t .

The a priori upper bound reads now as

Z
2

u
2 1
u
1 1 1

dt +
( )
+ ku( )k2H1 () = ku1 k2L2 () + ku0 k2H1 () , 0 ,

t (t)
2 2 t
2 0 2 2 0

0 L () L2 ()

and we can see that, due to the dissipation term on the left-hand side

Z

2

u (t)
dt > 0 and strictly increasing with ,

t
2

0 L ()

one has

d

E(t) < 0 ,

dt

i. e., the total energy of the system is not conserved.

8.5. EXERCISES 129

Time reversibility

Linear hyperbolicity allows us to reverse the time axis, and reconstruct the solution from knowing

its values at a final time, instead of at an initial time. In other words, the retrograde boundary-

value problem 2

w

w = 0 in Q ,

t2

w = 0 ) on (0, T ) , (8.4.2)

w = w 0

w on {T } ,

= w1

is well posed, as the forward problem. Indeed, appliying the change of variable = T t and

setting u(x , ) = w(x , T ), we have

u w

(x , ) = (x , T ) ,

t

2u 2w

(x , ) = (x , T ) ,

2 t2

u(x , ) = w(x , T )

and moreover

u(0) = w(0) = w0 ,

u w

(0) = (T ) = w1 ,

t

so that w satisfies the standard problem

2

u

u = f in Q ,

2

u = 0 ) on (0, T ) , (8.4.3)

u = w0

u on {0} .

= w1

Roughly speaking, we may say that w solves the problem (8.4.2) backwards in time.

This, however, no longer holds if we include the term wt , since in dissipative phenomena the

time axis cannot be reversed: indeed, the previous transformation would give the term u in

(8.4.3), with the wrong sign!

The fact the a dissipative retrograde problem is not well posed is nothing but the mathematical

counterpart of the physical concept of entropy.

8.5 Exercises

8.1. [... to be added ...]

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